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ACTA UNIVERSITATIS UPSALIENSIS<br />

Uppsala Dissertations in Mathematics<br />

39


Fredrik Strömberg<br />

<strong>Computational</strong> <strong>Aspects</strong> <strong>of</strong> <strong>Maass</strong><br />

Waveforms


Dissertation at Uppsala University to be publicly examined in Room 2347, Polacksbacken,<br />

Tuesday, March 15, 2005 at 13:15 for the Degree <strong>of</strong> Doctor <strong>of</strong> Philosophy. The examination<br />

will be conducted in English.<br />

Abstract<br />

Strömberg, F. 2005. <strong>Computational</strong> <strong>Aspects</strong> <strong>of</strong> <strong>Maass</strong> Waveforms. Uppsala Dissertations in<br />

Mathematics 39. x, 176 pp. Uppsala. ISBN 91-506-1794-X<br />

The topic <strong>of</strong> this thesis is computation <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong>, and we consider a number<br />

<strong>of</strong> different cases: Congruence subgroups <strong>of</strong> PSL(2,Z) and Dirichlet characters (chapter 1);<br />

congruence subgroups and general multiplier systems and real weight (chapter 2); and noncongruence<br />

subgroups <strong>of</strong> the Modular group (chapter 3). In each case we first discuss the necessary<br />

theoretical background. We then outline the algorithm and display some <strong>of</strong> the results obtained<br />

by it.<br />

Keywords: <strong>Maass</strong> waveform, Congruence group, Noncongruence group, Multiplier system,<br />

Theta multiplier system, Eta multiplier system, weight, Hecke operators, Involution, Dirichlet<br />

character, Spectral theory, Computation, Whittaker function<br />

Fredrik Strömberg, Department <strong>of</strong> Mathematics, Uppsala University, Lägerhyddsvägen 2 ,Box<br />

480, SE-75106 Uppsala, Sweden<br />

c○ Fredrik Strömberg 2005<br />

ISBN 91-506-1794-X<br />

ISSN 1401-2049<br />

urn:nbn:se:uu:diva-4778 (http://urn.kb.se/resolve?urn=urn:nbn:se:uu:diva-4778)


Contents<br />

Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 1<br />

1 <strong>Maass</strong> Waveforms on Hecke Congruence subgroups with Dirichlet<br />

characters<br />

(<strong>Computational</strong> <strong>Aspects</strong>) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 11<br />

1.1 General Definitions and Notation . . . . . . . . . . . . . . . . . . . . . . . 12<br />

1.1.1 A Brief Introduction to Fuchsian Groups . . . . . . . . . . . . . 12<br />

1.1.2 Hecke Congruence Groups . . . . . . . . . . . . . . . . . . . . . . . 13<br />

1.1.3 Introduction to Dirichlet Characters . . . . . . . . . . . . . . . . 16<br />

1.1.4 A Brief Introduction to <strong>Maass</strong> Waveforms . . . . . . . . . . . . 18<br />

1.2 Some Structural Theory <strong>of</strong> M(Γ0(N), χ,λ) . . . . . . . . . . . . . . . 20<br />

1.2.1 The Conjugation Operator . . . . . . . . . . . . . . . . . . . . . . . . 20<br />

1.2.2 The Fourier Series . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 20<br />

1.2.3 Involutions and Normalizers . . . . . . . . . . . . . . . . . . . . . . 21<br />

1.2.4 The Reflection Operator . . . . . . . . . . . . . . . . . . . . . . . . . 22<br />

1.2.5 Complete Symmetrization . . . . . . . . . . . . . . . . . . . . . . . 23<br />

1.2.6 Hecke Operators . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 25<br />

1.2.7 Oldforms and Newforms . . . . . . . . . . . . . . . . . . . . . . . . . 28<br />

1.2.8 The Cusp Normalizing Maps as Normalizers <strong>of</strong> (Γ0(N), χ) 32<br />

1.3 <strong>Computational</strong> <strong>Aspects</strong> . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 34<br />

1.3.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 34<br />

1.3.2 Phase 1 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 34<br />

1.3.3 Phase 2 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 43<br />

1.3.4 Remarks on the Performance <strong>of</strong> the Algorithm . . . . . . . . . 44<br />

1.4 Results . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 46<br />

1.4.1 Eigenvalues . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 47<br />

1.4.2 Lowest Eigenvalues . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 47<br />

1.4.3 Fourier Coefficients . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 48<br />

2 Computation <strong>of</strong> <strong>Maass</strong> Waveforms with Non-trivial Multiplier Systems<br />

. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 55<br />

2.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 55<br />

2.2 Multiplier Systems . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 55<br />

2.2.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 55<br />

2.2.2 The η multiplier system . . . . . . . . . . . . . . . . . . . . . . . . . . 57<br />

2.2.3 The θ Multiplier System . . . . . . . . . . . . . . . . . . . . . . . . . 59<br />

vii


2.2.4 Further properties <strong>of</strong> the multiplier systems . . . . . . . . . . . 61<br />

2.3 <strong>Maass</strong> Waveforms . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 62<br />

2.3.1 Decomposition <strong>of</strong> the discrete spectrum . . . . . . . . . . . . . . 64<br />

2.4 Operators . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 65<br />

2.4.1 Conjugation and reflection . . . . . . . . . . . . . . . . . . . . . . . . 65<br />

2.4.2 The involution ωN . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 68<br />

2.4.3 The operator σ4 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 70<br />

2.4.4 <strong>Maass</strong> operators . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 71<br />

2.4.5 <strong>Maass</strong> operators and the symmetry about k = 6 . . . . . . . . 72<br />

2.4.6 Hecke operators . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 73<br />

2.4.7 Lifts at weight 1 and Fourier coefficients . . . . . . . . . . . . . 77<br />

2.5 Oldforms . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 80<br />

2.5.1 Γ0(N) with N prime and η-multiplier . . . . . . . . . . . . . . . 81<br />

2.6 The Eisenstein series . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 82<br />

2.6.1 Weight 0 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 82<br />

2.7 Half integer weight . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 83<br />

2.7.1 The Shimura correspondence . . . . . . . . . . . . . . . . . . . . . . 83<br />

2.8 Some <strong>Computational</strong> Remarks . . . . . . . . . . . . . . . . . . . . . . . . . 86<br />

2.9 Numerical results . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 89<br />

2.9.1 Varying weight . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 89<br />

2.9.2 Small weights . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 90<br />

2.9.3 Lifts at weight 1 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 94<br />

2.9.4 Half integer weight . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 94<br />

3 Computation <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong> on Non-congruence Subgroups <strong>of</strong><br />

the Modular Group . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 107<br />

3.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 107<br />

3.2 Subgroups <strong>of</strong> the Modular Group . . . . . . . . . . . . . . . . . . . . . . . 108<br />

3.2.1 Permutations and subgroups <strong>of</strong> PSL(2,Z) . . . . . . . . . . . . 109<br />

3.2.2 Congruence subgroups . . . . . . . . . . . . . . . . . . . . . . . . . . . 118<br />

3.3 Numerical implementation . . . . . . . . . . . . . . . . . . . . . . . . . . . 119<br />

3.4 Numerical Results . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 121<br />

3.4.1 Subgroups . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 121<br />

3.4.2 Newforms . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 123<br />

3.4.3 Fourier coefficients . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 123<br />

4 An Algorithm for Whittaker’s W-function . . . . . . . . . . . . . . . . . . . 141<br />

4.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 141<br />

4.2 Presentation <strong>of</strong> the Algorithm . . . . . . . . . . . . . . . . . . . . . . . . . 141<br />

4.3 The path . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 142<br />

4.4 The integrand . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 147<br />

4.5 Evaluation <strong>of</strong> the Integrand . . . . . . . . . . . . . . . . . . . . . . . . . . . 150<br />

4.5.1 The power series . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 151<br />

4.5.2 Recursion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 153<br />

viii


4.5.3 Asymptotic series . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 155<br />

4.5.4 Comparison . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 155<br />

4.6 Evaluation <strong>of</strong> the Integrals . . . . . . . . . . . . . . . . . . . . . . . . . . . . 155<br />

4.7 Remarks on the K-Bessel function . . . . . . . . . . . . . . . . . . . . . . 156<br />

4.8 Numerical verification . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 157<br />

Swedish Summary . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 159<br />

Acknowledgements . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 167<br />

References . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 169<br />

ix


Introduction<br />

This is a dissertation in experimental mathematics, and we are thus in a situation<br />

where the “pure” mathematics takes advantage <strong>of</strong> numerical results to<br />

build hypotheses.<br />

A short description <strong>of</strong> what distinguishes experimental mathematics is that<br />

by using theoretical reasoning you construct numerical models (algorithms)<br />

which are then used to generate information. On the basis <strong>of</strong> the numerical<br />

results one then formulates hypotheses and conjectures (for a more detailed<br />

description see [16]).<br />

The prospect is <strong>of</strong> course that later one will be able provide rigorous pro<strong>of</strong>s<br />

<strong>of</strong> these conjectures, but even in those situations that the present-day mathematical<br />

machinery is not up to this, the numerical results can provide valuable<br />

insight and inspire to further studies.<br />

The three chapters <strong>of</strong> the dissertation are all designed in the same way; the<br />

bulk is devoted to a presentation <strong>of</strong> the theory, which is then used to construct<br />

the algorithms. After a brief description <strong>of</strong> how the algorithms are designed,<br />

we present our experimental findings.<br />

The problem we have studied is the construction <strong>of</strong> the mathematical objects<br />

referred to as “<strong>Maass</strong> <strong>waveforms</strong>” (after Hans <strong>Maass</strong> (1911-1992)). Mathematical<br />

applications <strong>of</strong> these wavefunctions are mainly in analytical number<br />

theory. There are also connections to physics, foremost to mathematical models<br />

<strong>of</strong> quantum chaos ([35, 14]), but also to cosmology ([8]).<br />

On the Physical Background<br />

Consider the motion <strong>of</strong> a free (quantum mechanical) particle with mass m0<br />

on a surface S . According to quantum mechanics, this system is described<br />

by a wave function, ψ, which tells us about the probability <strong>of</strong> encountering<br />

the particle in a certain region <strong>of</strong> the surface. The propagation <strong>of</strong> this wave<br />

function, for a particle with the energy E = ¯hv, is governed by the Schrödinger<br />

equation:<br />

i¯h ∂ψ<br />

∂t<br />

¯h2<br />

= − ∆ψ = Eψ,<br />

2m0<br />

1


where ∆ is the so-called Laplace-Beltrami operator on the surface. (e.g. in the<br />

plane: ∆ψ = ψ ′′<br />

xx + ψ ′′<br />

yy). We now assume that we can separate ψ as a product<br />

<strong>of</strong> a time-dependent part, e−ivt , and a time-independent part, φ(P), P ∈ S . We<br />

thus obtain an equation for the time independent part (the time-independent<br />

Schrödinger equation, or the Laplace equation):<br />

∆φ + λφ = 0, (⋆)<br />

where we have set λ = 2m0E<br />

¯h 2 (the eigenvalue), and since the probability is one<br />

that the particle is somewhere on the surface we must have<br />

<br />

|φ| 2 dP = 1.<br />

S<br />

If we have a charged particle and introduce a constant magnetic field, B, the<br />

time independent part <strong>of</strong> the wave function will obtain a “twist” (spin) and φ<br />

will now satisfy a modified Laplace equation:<br />

∆Bφ + λφ = 0, (⋆⋆)<br />

where ∆B is a modified Laplace-Beltrami operator, (in our case, it will read<br />

∆B = ∆ − iyB ∂<br />

∂x , with B proportional to the magnetic field strength).<br />

Observe that the wave equation (⋆) is “universal”, in the sense that it describes<br />

many different kinds <strong>of</strong> wave motion. For example, if S describes a<br />

drum, then (⋆) will describe the appearance <strong>of</strong> a standing wave on the drum<br />

skin, with a frequency proportional to √ λ. Bass notes, thus correspond to<br />

small λ and higher notes to larger λ.<br />

In this thesis we consider the equations (⋆) and (⋆⋆) for certain special surfaces,<br />

S , namely those with constant negative curvature (we say that the surface<br />

is hyperbolic), finite area and a finite (positive) number <strong>of</strong> “punctures”<br />

or “cusps” (i.e. we have removed a finite number <strong>of</strong> points from the surface<br />

and pulled them away towards infinity). Furthermore, the surfaces have to<br />

be special in an arithmetical way, which in this case means that they correspond<br />

to subgroups <strong>of</strong> the modular group. We also allow our surfaces to have<br />

“corners” (marked points). The solutions to (⋆) and (⋆⋆) on such surfaces are<br />

called <strong>Maass</strong> <strong>waveforms</strong>.<br />

One reason that surfaces <strong>of</strong> negative curvature are <strong>of</strong> interest in this context<br />

is that the classical free-particle motion on these surfaces is known to be<br />

(strongly) chaotic.<br />

By studying the solutions <strong>of</strong> (⋆) on such a surface we might thus gain some<br />

insight about “traces” <strong>of</strong> this classical chaos left behind in the quantized system.<br />

This type <strong>of</strong> investigations is part <strong>of</strong> the research area quantum chaos.<br />

The wavefunctions are not chaotic in themselves, and the classical chaos is<br />

instead visible in e.g. the distribution <strong>of</strong> the energy levels (eigenvalues) and<br />

2


in the appearance <strong>of</strong> the wavefunctions φ which corresponds to high energies.<br />

In Figure 2 you can see a sample <strong>of</strong> how a <strong>Maass</strong> waveform, ϕ with an<br />

eigenvalue λ ≈ 10000 looks. (a) Shows the nodal lines (i.e. the points where<br />

ϕ(z) = 0), and in (b) you can see the probability function associated to ϕ,<br />

i.e. |ϕ(z)| 2 . The light regions in (b) shows the places in the plane where the<br />

quantum mechanical particle likes to hang around.<br />

The main purpose with the work underlying this thesis was to obtain stable<br />

algorithms to compute eigenvalues, λ, and corresponding wave functions, φ,<br />

to the above mentioned problems.<br />

A brief Introduction to Hyperbolic Geometry<br />

To visualize negative curvature (hyperbolic geometry), as a mathematician<br />

one usually works in the upper half-plane model, i.e.<br />

H = {x + iy ∈ C|y > 0},<br />

where we introduce a concept <strong>of</strong> distance and area with the help <strong>of</strong> the arc<br />

element<br />

ds = |dz|<br />

y ,<br />

and the area element<br />

dµ = dxdy<br />

.<br />

y2 The shortest path between two points z1,z2 ∈ H is obtained by following<br />

the geodesic which passes through these points. In the usual (flat) case the<br />

geodesic is the unique straight line which passes through these points. In the<br />

hyperbolic upper half plane, the geodesics are the semi-circles and straight<br />

lines which are perpendicular to the x-axis. See Figure 1 for an illustration <strong>of</strong><br />

Figure 1: The upper half-plane<br />

H<br />

γ3<br />

y<br />

z4<br />

z3<br />

γ1<br />

γ2<br />

z2<br />

z1<br />

(a) Geodesics and a horocycle<br />

x<br />

S<br />

F1 F4 F7<br />

E<br />

e2<br />

e3 e ′ 3<br />

F6<br />

F3<br />

F5<br />

F2<br />

(b) A fundamental domain<br />

3


the geodesics γ1 and γ2 which connect z1 with z2 and z3 with z4 respectively.<br />

In the figure we also give an example <strong>of</strong> a horocycle, γ3.<br />

A consequence <strong>of</strong> this (odd) way to measure distances in the upper halfplane<br />

is for example that you can never get all the way down to the real line,<br />

y = 0 (i.e. the distance from any point in H to the real line is infinite).<br />

A more physical way to think about the hyperbolic space is the following<br />

(after Poincaré [82]):<br />

"Visualize yourself in a finite spherical universe, bounded by a “shell”. Assume<br />

that the temperature in this universe is proportional to the distance from<br />

the shell. That is, it is hottest in the center, and the absolute zero temperature is<br />

attained at the shell. Now, we also make the assumption that the speed <strong>of</strong> movement<br />

is proportional to the temperature (this is what one expect physically, at<br />

the molecular/atomic level).<br />

If we now decide to move to the end <strong>of</strong> the universe (i.e. the shell) it is easy<br />

to realize that we will never get there (in finite time). What we will experience<br />

is that the closer we get to the shell, the lower the temperature gets and the<br />

slower we move.<br />

Mathematically, under these assumptions, we can can compute our distance<br />

from the center at the time t, r(t) = R(1 − e −ct ), where R is the radius <strong>of</strong> the<br />

shell and c is a positive constant."<br />

An isometry in a space (where we have a concept <strong>of</strong> distances) is a mapping<br />

which preserves distances and areas. In the usual “flat” case, the translations,<br />

rotations and reflections are examples <strong>of</strong> isometries. Reflection is an example<br />

<strong>of</strong> an orientation-reversing isometry, while translations and rotations are examples<br />

<strong>of</strong> orientation-preserving isometries. In the upper half-plane we know<br />

that any orientation preserving isometry can be written as a Möbius transfor-<br />

mation:<br />

T : z ↦→<br />

az + b<br />

, where a,b,c,d ∈ R, ad − bc = 1.<br />

cz + d<br />

Such a transformation is then associated to a pair <strong>of</strong> matrices:<br />

<br />

±<br />

a<br />

c<br />

b<br />

d<br />

,<br />

and the group <strong>of</strong> all such isometries are usually denoted by PSL(2,R) =the<br />

group <strong>of</strong> all matrices <br />

a b<br />

c d , with a,b,c,d ∈ R, ad − bc = 1, and where we<br />

identify the matrices <br />

a b −a −b<br />

c d and −c −d .<br />

One knows that all hyperbolic surfaces can be described with the help <strong>of</strong><br />

discrete subgroups <strong>of</strong> PSL(2,R), i.e. we can write our surface S as S =<br />

Γ\H , where Γ is a Fuchsian group (a discrete group <strong>of</strong> Möbius transformations).<br />

Those surfaces, S , we are interested in have “horns” (points at ∞) and<br />

4


they can also have corners. In the corresponding group, Γ, the horns show<br />

themselves as fixed-points to parabolic transformations, i.e. those that correspond<br />

to matrices with the absolute value <strong>of</strong> the trace = 2 (remember that the<br />

trace <strong>of</strong> a matrix <br />

a b<br />

c d is defined as a + d). The corners, on the other hand are<br />

fixed-points <strong>of</strong> elliptic (|trace|< 2). Therefore, one usually says that the corresponding<br />

points are parabolic or elliptic. (The parabolic fixed-points are also<br />

called cusps). In PSL(2,R), in addition to elliptic and parabolic transformations<br />

(which geometrically correspond to rotations and translations) also hyperbolic<br />

transformations (|trace|> 2). While the first types have unique fixedpoints,<br />

in general the hyperbolic maps have two fixed points in R ∪ {∞} (and<br />

geometrically, they describe a movement along the geodesic which connects<br />

these points).<br />

To illustrate the surface S = Γ\H it is customary to use a fundamental<br />

domain, which in our case can be taken as a connected domain F ⊂ H with<br />

the following two properties: we can not map any point inside F to another<br />

point inside F with a map in Γ, and secondly, the Γ-images F must tessellate<br />

H (i.e. ∪γ∈Γγ(F ) = H). An example <strong>of</strong> a fundamental domain is the<br />

standard fundamental domain for the modular group, PSL(2,Z) (all matrices<br />

in PSL(2,Z) with integer entries), F1 = {z = x + iy ∈ H||z| ≥ 1, |x| ≤ 1<br />

2 }.<br />

In Figure 1(b) we can see F1 together with 6 “copies”. Observe that all images,<br />

F j, have the same hyperbolic area. If we identify the sides <strong>of</strong> F1 using<br />

the maps S : z ↦→ z + 1 and E : z ↦→ −1 we get (topologically) a sphere with<br />

z<br />

one cusp (i∞) and two corners (e2 and e3 which is identified with e ′ 3<br />

). One can<br />

show that the side pairing transformations E and S generates PSL(2,Z) and,<br />

in fact, in all cases we consider here, we can always choose the fundamental<br />

domain in such a way that the side pairing maps generate the group.<br />

Summary <strong>of</strong> the Chapters in the Thesis<br />

In chapter one we study the problem which in the upper half-plane model can<br />

be written as:<br />

y 2 φ ′′<br />

xx(z) + φ ′′<br />

yy(z) + λφ(z) = 0, z = x + iy ∈ H , (1)<br />

φ (T z) = χ(T )φ(z),<br />

<br />

T ∈ Γ, z ∈ H , (2)<br />

|φ| 2 dµ = 1, (3)<br />

Γ\H<br />

where χ : Γ → S 1 as an even Dirichlet character (i.e. χ(AB) = χ(A)χ(B) and<br />

χ(−1) = 1), and Γ is a Hecke-congruence subgroup <strong>of</strong> the modular group. For<br />

a given positive integer N, we define the Hecke-congruence group, Γ0(N), in<br />

5


the following way:<br />

<br />

<br />

Γ0(N) =<br />

a<br />

c<br />

b<br />

d<br />

|a,b,c,d ∈ Z, ad − bc = 1, c ≡ 0 mod N .<br />

We will now briefly outline the algorithm which we use to find solutions <strong>of</strong> this<br />

problem. Given an eigenvalue, λ = 1<br />

4 + R2 , we know that the corresponding<br />

φ can be written as a Fourier series:<br />

φ(x + iy) = ∑ c(n)<br />

n=0<br />

√ yKiR(2π|n|y)e 2πinx ,<br />

where KiR(y) is the K-Bessel function. Due to the rapid decay <strong>of</strong> KiR(y) we<br />

can approximate φ (up to any given ε > 0) with a truncated Fourier series,<br />

φ(x + iy) ≈ ˆφ(x + iy) = ∑<br />

0 0 is chosen wisely), thus providing us with a set <strong>of</strong> linear equations for<br />

the coefficients c(n):<br />

c(n) √ Y KiR(2π|n|Y ) = 1<br />

2Q<br />

Q<br />

∑<br />

m=1−Q<br />

ˆφ(zm)e(−nxm) ≈ 1<br />

2Q<br />

Q<br />

∑<br />

m=1−Q<br />

φ(zm)e(−nxm),<br />

and by using the implicit automorphy (2) we can replace the point zm with a<br />

point, z ∗ m = x ∗ m + iy ∗ m in the fundamental domain <strong>of</strong> the group. Doing this, we<br />

get a non-trivial, well-conditioned, linear system <strong>of</strong> equations<br />

c(n) √ Y KiR(2π|n|Y ) = ∑ Vnkc(k),<br />

|k|≤M0<br />

(4)<br />

which we can solve to obtain the Fourier coefficients c(n).<br />

If we do this for an arbitrary value <strong>of</strong> R, the resulting function φ will obviously<br />

not be invariant under Γ. But if 1<br />

4 + R2 is close to a true eigenvalue<br />

the computed Fourier coefficients will be close to the coefficients <strong>of</strong> a true<br />

eigenfunction. In particular, for a “true” R, the solution vectors should be independent<br />

<strong>of</strong> the value <strong>of</strong> Y (provided it is kept small enough that the system<br />

does not become trivial).<br />

We use this fact to locate eigenvalues as follows; using two different values<br />

<strong>of</strong> Y in parallel, we solve (4) over and over again, for R running through a<br />

specified interval, all the while we try to minimize the functional<br />

6<br />

h(R) = c(2) − c ′ (2) + c(3) − c ′ (3) + c(4) − c ′ (4) ,


where c(n) and c ′ (n) are the solutions <strong>of</strong> (4) corresponding to the different<br />

Y ’s. When we find a minimizing R-value, we go through with further (more<br />

technical, e.g. using Hecke-relations) tests to see if we indeed seem to be close<br />

to a true eigenvalue.<br />

As soon as we have an approximate eigenvalue and a corresponding set <strong>of</strong><br />

Fourier coefficients we can use another algorithm (phase 2) to easily compute<br />

many more coefficients.<br />

In chapter 2 we study a similar problem, where we instead <strong>of</strong> (1) consider<br />

the equation<br />

y 2 φ ′′<br />

xx(z) + φ ′′<br />

yy(z) − iymφ ′ x + λφ(z) = 0,z = x + iy ∈ H , (1’)<br />

where the weight, m, is a real number, and where the condition (2) is replaced<br />

by<br />

φ(T z) = v(T ) jT (z;m)φ(z), (2’)<br />

where v : Γ → C is a so-called multiplier system and<br />

jT (z;m) = exp(2πimArg(cz + d)),<br />

for T = <br />

a b<br />

c d and m ∈ R.<br />

Physically, we can interpret the equation (1’) as if we now have a charged<br />

(quantum mechanical) particle moving in a constant magnetic field on the surface<br />

Γ\H. The magnetic field strength is proportional to m, and the multiplier<br />

system can be used to describe the magnetic flux through different parts <strong>of</strong> the<br />

surface (the horns, handles etc.) Cf. [2, 1, 11]. (For those that are knowledgeable<br />

in the physics here, (1’) corresponds to the choice <strong>of</strong> the Landau gauge,<br />

A = c m<br />

y dx for some suitable constant c.)<br />

It turns out that the method introduced in chapter one is still applicable after<br />

a few (minor) changes. The (computationally) hardest problem here is that the<br />

Fourier series now contain the Whittaker W-function instead <strong>of</strong> the K-Bessel<br />

function. The W-function is more difficult to compute, and we were forced<br />

to develop a new algorithm for that purpose. The algorithm is presented in<br />

Chapter 4.<br />

In chapter 3 we investigate the same problem as in chapter 1, but here<br />

we also consider subgroups Γ ⊆ PSL(2,Z) which are not congruence subgroups.<br />

The main difference between congruence subgroups and general noncongruence<br />

subgroups is that one doesn’t have as many “symmetries” available<br />

in the noncongruence case. This implies for example that one doesn’t<br />

have much a priori knowledge <strong>of</strong> neither the spectra nor the eigenfunctions.<br />

And one can therefore pose questions which are more or less answered in the<br />

congruence case: Are there any “new” <strong>Maass</strong> <strong>waveforms</strong>? If so, how many<br />

are there? Do they show any symmetries? How are their Fourier coefficients<br />

7


distributed?<br />

We prove the existence <strong>of</strong> new <strong>Maass</strong> <strong>waveforms</strong> for groups Γ which have<br />

certain specific properties, which on the surfaces Γ\H basically correspond to<br />

different forms <strong>of</strong> mirror symmetries. (To systematically verify these properties<br />

for large groups, computer aid is a necessity).<br />

8


Figure 2: Plots <strong>of</strong> a <strong>Maass</strong> waveform, ϕ(z) on Γ0(5) with R = 100.0870...<br />

2.1<br />

Nodal lines for R=100.0870... on Γ 0 (5)<br />

0.1<br />

−1 1<br />

(a) Nodal lines<br />

(b) Density plot, i.e. |ϕ(z)| 2<br />

9


1 <strong>Maass</strong> Waveforms on Hecke Congruence<br />

subgroups with Dirichlet characters<br />

(<strong>Computational</strong> <strong>Aspects</strong>)<br />

Introduction<br />

The main topic <strong>of</strong> this chapter is computational aspects <strong>of</strong> the theory <strong>of</strong> <strong>Maass</strong><br />

<strong>waveforms</strong>, i.e. square-integrable eigenfunctions <strong>of</strong> the Laplace- Beltrami operator,<br />

on certain Riemann surfaces with constant negative curvature and finite<br />

area. The surfaces under consideration correspond to quotients <strong>of</strong> the upper<br />

half-plane by certain discrete groups <strong>of</strong> isometries, the so called Hecke congruence<br />

subgroups.<br />

It is known that such functions can also be regarded as wavefunctions corresponding<br />

to a quantum-mechanical system describing a particle moving freely<br />

on the surface. Since the classical counterpart <strong>of</strong> this motion is chaotic, the<br />

study <strong>of</strong> these wavefunctions is closely related to the study <strong>of</strong> quantum chaos<br />

on the surface. The presentation here, however, will be purely from a mathematical<br />

viewpoint.<br />

Today, our best knowledge <strong>of</strong> generic <strong>Maass</strong> <strong>waveforms</strong> comes from numerical<br />

experiments, and these have previously been limited to the modular<br />

group, PSL(2,Z), and certain triangle groups (cf. [43, 45, 42] and [98]). One<br />

<strong>of</strong> the primary goals <strong>of</strong> this lecture is to give the necessary theoretical background<br />

to generalize these numerical experiments to Hecke congruence subgroups<br />

and non-trivial characters. We will describe algorithms that can be<br />

used to locate eigenvalues and eigenfunctions, and we will also present some<br />

<strong>of</strong> the results obtained by those algorithms.<br />

There are four sections, first elementary notations and definitions, mostly<br />

from the study <strong>of</strong> Fuchsian groups and hyperbolic geometry. Then more <strong>of</strong> the<br />

theoretical background needed to understand the rich structure <strong>of</strong> the space <strong>of</strong><br />

<strong>Maass</strong> <strong>waveforms</strong> will be introduced. The third section deals with the computational<br />

aspects, and the final section contains some numerical results.<br />

11


1.1 General Definitions and Notation<br />

1.1.1 A Brief Introduction to Fuchsian Groups<br />

For a more thorough (but still elementary) introduction to this subject see for<br />

example [52] or [31].<br />

The following notation will be used; M2(K) is the ring <strong>of</strong> 2 × 2 matrices<br />

over some base ring K, usually R or Z, GL(2,K) ⊆ M2(K) is the group <strong>of</strong><br />

invertible 2 × 2 matrices and SL(2,K) ⊆ M2(K) is the group <strong>of</strong> matrices with<br />

determinant equal to 1. H denotes the Poincaré upper half-plane {z = x +<br />

iy|y > 0} equipped with the hyperbolic metric and area measure<br />

ds 2 = 1<br />

y (dx2 + dy 2 ), dµ = 1<br />

y<br />

2 dxdy.<br />

PSL(2,R) is the group <strong>of</strong> Möbius transformations with real coefficients,<br />

<br />

<br />

az + b <br />

PSL(2,R) = z → <br />

cz + d a,b,c,d ∈ R, ad − bc = 1 .<br />

It is clear that we can represent such mappings with matrices and there is an<br />

isomorphism PSL(2,R) ≈ SL(2,R)/{±Id}, where Id is the identity element<br />

in SL(2,R). We will use matrix and transformation notation interchangeably,<br />

all groups discussed are subgroups <strong>of</strong> PSL(2,R) so they always contain −Id,<br />

hence there should be no confusion when we use the same notation for matrix<br />

groups as for transformation groups. Note that H, the hyperbolic metric and<br />

the hyperbolic measure are invariant under PSL(2,R). The only subgroups <strong>of</strong><br />

PSL(2,R) that we are interested in are the discrete subgroups.<br />

Definition 1.1.1. A Fuchsian group is a discrete subgroup <strong>of</strong> PSL(2,R).<br />

The basic example <strong>of</strong> a Fuchsian group, for us, is the modular group,<br />

PSL(2,Z) = SL(2,Z)/{±Id} ⊆ PSL(2,R).<br />

This is just the subgroup <strong>of</strong> Möbius transformations in PSL(2,R) with integer<br />

coefficients. This is also an example <strong>of</strong> what is called an arithmetic group<br />

(cf. [52, ch. 5]) and is <strong>of</strong> interest for number theorists among others.<br />

If Γ is a Fuchsian group it is a standard result that the space <strong>of</strong> Γ−orbits,<br />

Γ\H = {Γz|z ∈ H}, can be given the analytical structure <strong>of</strong> a Riemann surface<br />

with marked points (this is also called an orbifold) (cf. [65, II.F.], and for<br />

more details [85]). On the other hand, a classical result, the Klein-Poincaré<br />

uniformization theorem, asserts that any Riemann surface M with constant<br />

negative curvature equal to −1 can be realized as Γ\H for some Fuchsian<br />

group Γ (see [60] for more details on uniformization). We can visualize the<br />

Riemann surface Γ\H via a fundamental domain F for Γ.<br />

12


Definition 1.1.2. A closed set F ⊆ H is a fundamental domain for Γ if<br />

a) ∪T ∈ΓT (F) = H, and<br />

b) if F ◦ denotes the interior <strong>of</strong> F then T1(F ◦ ) ∩ T2(F ◦ ) = /0 if T1 =<br />

T2 ∈ Γ.<br />

Note that we use closed fundamental domains in line with for example [52],<br />

so they will in general contain some equivalent boundary points. There are<br />

two important examples <strong>of</strong> fundamental domains.<br />

Definition 1.1.3. Let Γ be a Fuchsian group and let p ∈ H be a point that<br />

is not fixed by any element in Γ. Let d(z,w) denote the hyperbolic distance<br />

between z and w and define<br />

DΓ(p) = {z ∈ H|d(z, p) ≤ d(z,T p), ∀T ∈ Γ}.<br />

The set DΓ(p) is called a Dirichlet fundamental domain for Γ, with center p.<br />

Let Γz = {T ∈ Γ|T z = z} denote the stabilizer in Γ <strong>of</strong> the point z ∈ H ∪R∪<br />

{∞}. Another important type <strong>of</strong> fundamental domain, FΓ, the Ford fundamental<br />

domain, can be described in terms <strong>of</strong> the exterior <strong>of</strong> isometric circles. Let<br />

DT = z ∈ H |T ′ (z)| ≥ 1 ,∀T ∈ PSL(2,R). Then, if Γ∞ = {T ∈ Γ|T ∞ = ∞}<br />

is non-empty and generated by S : z ↦→ z + 1, we define<br />

FΓ = <br />

T ∈Γ\Γ∞<br />

DT ∩<br />

<br />

z ∈ H |ℜz| ≤ 1<br />

<br />

.<br />

2<br />

The boundary (in H ∪ R ∪ {∞}) <strong>of</strong> a fundamental domain F might contain<br />

vertices that are fixed-points <strong>of</strong> elements <strong>of</strong> Γ; these points are called elliptic<br />

or parabolic vertices respectively if the transformations fixing them are either<br />

elliptic or parabolic. A parabolic vertex is usually referred to as a cusp, and is<br />

viewed as either a point removed from the surface Γ\H, or a point at infinity.<br />

1.1.2 Hecke Congruence Groups<br />

Let N be any positive integer. We define the principal congruence subgroup<br />

<strong>of</strong> level N, Γ(N) ⊆ PSL(2,Z), by<br />

<br />

<br />

a b<br />

<br />

Γ(N) =<br />

∈ PSL(2,Z) <br />

a b 1 0<br />

c d<br />

≡<br />

mod N .<br />

c d 0 1<br />

This is a subgroup <strong>of</strong> finite index in PSL(2,Z), and any subgroup <strong>of</strong> PSL(2,Z)<br />

containing some Γ(N) is called a congruence subgroup. The Hecke congru-<br />

13


ence subgroup, Γ0(N), is defined by<br />

<br />

<br />

<br />

a b<br />

<br />

Γ0(N) =<br />

∈ PSL(2,Z) <br />

c d<br />

c ≡ 0 mod N .<br />

It is obvious that Γ(N) ⊆ Γ0(N) so Γ0(N) is a congruence subgroup.<br />

The standard fundamental domain for PSL(2,Z) = Γ0(1) is the set<br />

<br />

F1 = z ∈ H <br />

<br />

1<br />

|z| ≥ 1, |ℜz| ≤ ,<br />

2<br />

which is also a Ford fundamental domain. Suppose from now on that we have<br />

fixed a set <strong>of</strong> right coset representatives {Vk} vN<br />

k=1 for Γ0(N) in Γ0(1). (See<br />

pp. 43–43 on how to construct such maps Vk.) One knows that<br />

−1<br />

vN = [Γ0(1) : Γ0(N)] = N∏ 1 + p<br />

p|N<br />

.<br />

We can now fix a fundamental domain for Γ0(N) corresponding to these coset<br />

representatives (cf. [52, thm. 3.1.2])<br />

FN = ∪ vN<br />

k=1 Vk(F1). (1.1)<br />

Note that FN need not be normal in the sense <strong>of</strong> [39], but it will be bounded<br />

by finitely many geodesics and it also has some other nice properties as we<br />

will see (Theorem 1.1). Denote by κ0 the total number <strong>of</strong> parabolic vertices<br />

<strong>of</strong> FN, and by κ the number <strong>of</strong> inequivalent cusps <strong>of</strong> FN (note κ ≤ κ0). It is<br />

known that<br />

<br />

κ = ∑ φ d,<br />

d|N,d>0<br />

N<br />

<br />

,<br />

d<br />

where φ is Eulers totient function,<br />

φ(n) = # j ∈ {1,2,3,...,n} ( j,n) = 1 <br />

(see [51, 89]). Let v1,...,vκ0 be the set <strong>of</strong> parabolic vertices <strong>of</strong> FN and<br />

p1,..., pκ a set <strong>of</strong> inequivalent cusps. Without loss <strong>of</strong> generality we set p1 =<br />

∞.<br />

Definition 1.1.4. A cusp normalizing map associated to a cusp p j, is a map<br />

σ j ∈ PSL(2,R) satisfying<br />

i) σ j(∞) = p j, and<br />

ii) σ jSσ −1<br />

j = S j,<br />

where Sz = z + 1 is a generator <strong>of</strong> Γ∞, and S j is a generator <strong>of</strong> Γp j , the (cyclic<br />

14


infinite) stabilizer group <strong>of</strong> the cusp p j. Then σ j is uniquely determined up to<br />

a translation.<br />

For the groups Γ0(N) we will always choose the cusp normalizing maps σ j<br />

<strong>of</strong> the form σ j = A jρ j, where A j ∈ SL(2,Z) maps ∞ to p j and ρ j is a scaling<br />

by h j, the so-called width <strong>of</strong> the cusp p j (cf. [51, pp. 36-37]). It is important<br />

to note that h j|N. For the cusp p1 = ∞ we will always take σ1 as the identity.<br />

We will now provide an important fact about the fundamental domain for<br />

Γ0(N) that will be <strong>of</strong> use to us later. First, for each parabolic vertex vℓ <strong>of</strong><br />

FN we choose a map Uℓ ∈ Γ0(N) which maps vℓ to its cusp representative in<br />

{p1,..., pκ}; say Uℓ(νℓ) = p j(ℓ). If vℓ is already a cusp representative we take<br />

Uℓ = Id. This will give us a finite collection <strong>of</strong> maps U = {Uℓ}1≤ℓ≤κ0 . Now,<br />

to any point w ∈ FN we associate the “closest” (see Remark 1.1.1) parabolic<br />

vertex, vw, and a corresponding map Uw = Uℓ(vw) ∈ U .<br />

Remark 1.1.1. By the “closest” vertex to a point w ∈ FN we mean the vertex<br />

vℓ with respect to which the point w has the greatest height, i.e. that vℓ for<br />

which ℑ(σ −1<br />

j(ℓ) Uℓw) is maximal. We use I(w) to denote the index <strong>of</strong> the cusp<br />

representative corresponding to the vertex closest to w.<br />

Define the height <strong>of</strong> w ∈ FN, h(w), by<br />

h(w) = ℑ<br />

<br />

σ −1<br />

I(w) Uw (w)<br />

and the minimal height <strong>of</strong> the fundamental domain FN, Y0, by<br />

<br />

Y0 = Y0(N) = inf h(w) = inf ℑ σ<br />

w∈FN w∈FN<br />

−1<br />

I(w) Uw<br />

<br />

(w) . (1.2)<br />

It is clear that Y0 depends only on the fundamental domain FN (i.e. the choice<br />

<strong>of</strong> the coset representatives {Vk}) and not on the choice <strong>of</strong> the cusp representatives<br />

{p j}. A standard compactness argument also implies that Y0 is strictly<br />

positive and larger than some fixed quantity only depending on the fundamental<br />

domain, and the following theorem actually gives us an explicit bound.<br />

Theorem 1.1. For any positive integer N, the minimal height <strong>of</strong> FN satisfies<br />

the following inequality<br />

√<br />

3<br />

Y0(N) ≥<br />

2N .<br />

Pro<strong>of</strong>. Let {Vk} be the fixed set <strong>of</strong> coset representatives <strong>of</strong> Γ0(N) in Γ0(1) and<br />

FN as above.<br />

<br />

Let w ∈ FN. We want to show that ℑ<br />

σ −1<br />

to Remark 1.1.1 it is enough to show that ℑ<br />

<br />

,<br />

<br />

I(w) Uww<br />

√<br />

3<br />

≥ 2N<br />

<br />

σ −1<br />

j(ℓ) Uℓw<br />

<br />

≥<br />

, and according<br />

√ 3<br />

2N<br />

for some ℓ.<br />

Since w ∈ FN = ∪Vk(F1) there is an index k such that w ∈ Vk(F1) = Rk. Let<br />

15


v ℓ(k) = Vk(∞) be the vertex at infinity <strong>of</strong> Rk, and p j(k) the cusp representative<br />

equivalent to v ℓ(k), i.e. U ℓ(k)v ℓ(k) = p j(k).<br />

Note that we can write σ −1<br />

j = l js jω j where l j is a translation, ω j ∈ PSL(2,Z)<br />

is such that ω jS jω −1<br />

j = Sh j , (where S h z = z + h) and s j is a scaling with 1<br />

h j ,<br />

where the positive integer h j is called the width <strong>of</strong> the cusp p j (cf. [51, p.<br />

36-37]). The width, h j is independent <strong>of</strong> the choice <strong>of</strong> cusp-representative in<br />

the class <strong>of</strong> p j (cf. [63, p. 59]). Now ω j(k)U ℓ(k)Vk(F1) is simply a horizontal<br />

translate <strong>of</strong> F1, so ℑ(ω j(k)U ℓ(k)w) ≥<br />

√ 3<br />

2<br />

and thus ℑ(σ −1<br />

j(k) U ℓ(k)w) ≥<br />

[51, p. 36-37] it is clear that 1 ≤ h j ≤ N for all j, hence ℑ(σ −1<br />

and accordingly Y0 ≥<br />

√ 3<br />

2N .<br />

√<br />

3<br />

2h . From<br />

j(k)<br />

j(k) U √<br />

3<br />

ℓ(k)w) ≥ 2N ,<br />

Remark 1.1.2. If P is a prime ≥ 3 we will see that for a natural choice (which √<br />

3<br />

we use in the numerical work) <strong>of</strong> fundamental domain we have Y0(P) = 2P .<br />

Let T z = −1<br />

z and Sz = z + 1. We know that the maps V1 = Id and Vk =<br />

P+3<br />

T S (k− 2 ) for k = 2,...P + 1 is a set <strong>of</strong> right coset representatives for Γ0(P)<br />

in Γ0(1) (cf. [83, p. 135]). The corresponding fundamental domain is FP =<br />

∪ P+1<br />

k=1Vk(F1), it has one cusp <strong>of</strong> width 1 at p1 = ∞ and one cusp <strong>of</strong> width P at<br />

p2 = 0. The cusp normalizing maps are σ1 = Id and σ2z = −1<br />

1<br />

Pz . Let e = 2 +i<br />

√<br />

3<br />

2 .<br />

√ <br />

3 <br />

Then we see that ℑ(VP+1(e)) = 2 1<br />

√ <br />

3 P−1<br />

2 + i 2 + 2 −2 √<br />

3 4 = 2 P2 +3 ≤<br />

√<br />

3<br />

2P (since<br />

P ≥ 3), and clearly ℑ(σ −1<br />

2 VP+1(e)) = ℑ √<br />

e 1 1 3<br />

P + 2 − 2P = 2P . Hence the height<br />

√<br />

3<br />

<strong>of</strong> VP+1(e) is equal to 2P , which means that Y0<br />

√<br />

3<br />

≤ 2P . The conclusion is that<br />

for a prime level P ≥ 3 and for this particular choice <strong>of</strong> fundamental domain<br />

we have Y0 =<br />

√ 3<br />

2P .<br />

1.1.3 Introduction to Dirichlet Characters<br />

First we fix m ∈ Z + . A Dirichlet character, χ mod m, is a group homomorphism<br />

from (Z/mZ) ∗ to the unit circle S 1 , which is viewed as a function on Z<br />

(with period m) by assigning χ(n) = 0 if (n,m) > 1. We then define the map<br />

χ : SL(2,Z) → S 1 by<br />

χ<br />

<br />

a b<br />

c d<br />

<br />

= χ(d).<br />

Note that χ = 0 on Γ0(m) since ad − bc = 1 and c ≡ 0 mod m imply that<br />

χ(ad) = χ(a)χ(d) = 1. Observe also that χ is a group homomorphism from<br />

Γ0(m) to S 1 .<br />

If χ has a period less than m (for values <strong>of</strong> n restricted by (n,m) = 1) then<br />

χ is said to be imprimitive, otherwise χ is said to be a primitive character.<br />

16


(Note that the trivial character, χ(n) ≡ 1, is imprimitive for m ≥ 2). If q is the<br />

smallest period <strong>of</strong> χ, then q is called the conductor <strong>of</strong> χ.<br />

If χ has period m and conductor q there is a unique way to define a character<br />

χ ′ mod m ′ , for all modulus m ′ which are multiples <strong>of</strong> q or multiples <strong>of</strong> m.<br />

We say that the character χ is even if χ(−1) = χ(1) = 1, and odd if χ(−1) =<br />

−χ(1) = −1.<br />

Since we are actually concerned with PSL(2,Z) we need χ(A) = χ(−A) to<br />

hold for any A ∈ SL(2,Z), hence we will only consider even characters.<br />

It is known (see [27]) that primitive real characters exist only for moduli <strong>of</strong><br />

the following types:<br />

m = N1,4N2, or 4N3,<br />

where N1 ≡ 1 mod 4, N2 ≡ 2 mod 4 and N3 ≡ 3 mod 4 and N1,N2 and N3 are<br />

square-free.<br />

One should note that these are precisely the fundamental discriminants <strong>of</strong><br />

real quadratic fields, and it is also a fact that all real characters are given by<br />

quadratic residue symbols (Kronecker’s extension <strong>of</strong> the Legendre symbol)<br />

cf., e.g. [25, p. 37].<br />

Since characters are multiplicative it is fairly easy to evaluate them at products<br />

<strong>of</strong> integers, but when it comes to linear combinations <strong>of</strong> integers it is<br />

harder. A useful trick for evaluating characters at certain linear combinations<br />

is to factor the character instead <strong>of</strong> the integer. Here we will give the simplest<br />

case, which will be <strong>of</strong> use for the pro<strong>of</strong> <strong>of</strong> Proposition 1.2.6 in Subsection<br />

1.2.8.<br />

Fact 1.1.1. Suppose that χ is a real primitive character mod N and that m ∈<br />

Z + is such that m|N and m, N<br />

<br />

m = 1. Then it is known (cf. [69, (3.1.4)]) that<br />

we can factor χ as χ = χmχ N , where χm and χ N are real primitive characters<br />

m<br />

m<br />

mod m and N<br />

m respectively.<br />

Example 1.1.1. Let χ, N and m be as above, and let t = qA + N<br />

q B for some<br />

A,B ∈ Z. Then<br />

<br />

χ(t) = χm(t)χ N (t) = χm mA +<br />

m<br />

N<br />

m B<br />

<br />

χ N mA +<br />

m<br />

N<br />

m B<br />

<br />

<br />

N<br />

= χm<br />

m B<br />

<br />

χ N (mA).<br />

m<br />

Definition 1.1.5. A character χ defined on Γ0(N) is said to leave the cusp<br />

p j open if it is trivial on the stabilizer subgroup Γp j , i.e. if χ(S j) = 1. If a<br />

character leaves all cusps open we say that the character is regular for Γ0(N).<br />

If a character does not leave the cusp open we say that the cusp is closed,<br />

and it is known that there is no contribution to the continuous part <strong>of</strong> the spectrum<br />

from a closed cusp (cf. [39, pp. 91-99]).<br />

17


Observe that the cusp at ∞ is fixed by S : z ↦→ z + 1 so χ(S) = 1 for any<br />

Dirichlet character. Since we know that S j = σ jSσ −1<br />

j we will also have<br />

χ(S j) = 1 whenever σ j is a (Γ0(N), χ)−normalizer as described in Definition<br />

1.2.1 below. In particular, for square free N ≡ 1 mod 4 all Dirichlet characters<br />

are regular (cf. [12]).<br />

The basic assumption from now on (unless anything else is explicitly stated)<br />

is that any character χ is a real and even Dirichlet character.<br />

1.1.4 A Brief Introduction to <strong>Maass</strong> Waveforms<br />

For a Fuchsian group Γ, a function f defined on the upper half-plane is called<br />

Γ−automorphic if it is invariant or transforms in a chosen manner (i.e. with<br />

respect to characters or general multipliers) under the action <strong>of</strong> Γ by the so-<br />

called slash operator <strong>of</strong> weight k ∈ Z defined by<br />

f |[k,A](z) =<br />

cz + d<br />

|cz + d|<br />

−k<br />

f (Az), A ∈ Γ.<br />

Such functions can be viewed as “living” on the Riemann surface Γ\H (or<br />

coverings <strong>of</strong> it in the case <strong>of</strong> characters). In this chapter we will only consider<br />

certain non-holomorphic L 2 −functions, the so called <strong>Maass</strong> <strong>waveforms</strong> (to be<br />

defined below) and only at weight k = 0 (note that f |[0,A](z) = f (Az)).<br />

In the hyperbolic metric on H the Laplace-Beltrami operator, ∆, takes the<br />

form<br />

∆ = y 2<br />

<br />

∂ 2 ∂ 2<br />

+<br />

∂x2 ∂y2 <br />

.<br />

Definition 1.1.6. If Γ is a Fuchsian group, a <strong>Maass</strong> waveform for Γ is a function<br />

f defined on H with the following properties<br />

i) ∆ f + λ f = 0, λ ≥ 0;<br />

ii) f (Az) = f (z), for all A ∈ Γ;<br />

iii)<br />

<br />

Γ\H | f |2 dµ < ∞.<br />

We usually write the eigenvalue as λ = 1<br />

4 + R2 . It is conjectured (Selberg)<br />

that for congruence subgroups, if λ = 0, then we can take R ∈ [0,∞) (λ ≥<br />

1<br />

1<br />

4 ); any eigenvalues λ ∈ (0, 4 ] are usually called exceptional. We denote the<br />

space <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong> for Γ by M(Γ), and the space corresponding to a<br />

particular eigenvalue by M(Γ,λ). If χ is a group homomorphism from Γ to<br />

S1 (cf. Subsection 1.1.3) we can replace condition ii) by<br />

ii’) f (Az) = χ(A) f (z), for all A ∈ Γ.<br />

18


We then say that f is a <strong>Maass</strong> waveform for (Γ, χ), the space <strong>of</strong> which we<br />

denote by M(Γ, χ) or M(Γ, χ,λ). It is also known that each space M(Γ, χ,λ)<br />

is finite dimensional (see for example [39, p. 140, thm. 11.11]), that <strong>Maass</strong><br />

<strong>waveforms</strong> span the discrete part <strong>of</strong> the spectra <strong>of</strong> ∆, and that the continuous<br />

part is spanned by the Eisenstein series (see for example [50, 97] or [39]).<br />

It is an important fact that if Γ is a congruence subgroup then every <strong>Maass</strong><br />

waveform f in M(Γ, χ,λ) with λ > 0 is a cusp form, meaning that it tends<br />

rapidly to 0 in each cusp. This fact follows from [39, pp. 327(note 15),<br />

78(claim 9.6), 284(line 1), 328(line -2)] since the scattering matrix <strong>of</strong> Γ(N)<br />

can be computed explicitly and shown not to have any poles for 1<br />

2 < s < 1; cf.,<br />

e.g., [75] or [48], and [69, p. 114].<br />

We can equip the space M(Γ, χ,λ) with the Petersson inner-product,<br />

<br />

< f ,g >= f ¯gdµ,<br />

where the integration can be taken over any fundamental domain for Γ. With<br />

this inner-product M(Γ, χ,λ) is now a finite dimensional Hilbert space.<br />

If Γ is a congruence subgroup and χ is trivial on Γ(N), the eigenvalues 0 =<br />

λ0 < λ1 ≤ λ2 ≤ ···, counted with multiplicity, form a discrete sequence and<br />

a general Weyl’s law is known to hold. In particular, for Hecke congruence<br />

subgroups and the trivial character, we have (see [86, thm. 2])<br />

√ <br />

T<br />

N Γ0(N)(T ) = µ(FN)<br />

4π<br />

as T → ∞,<br />

Γ\H<br />

2κ √ √ √<br />

T − T ln T + A T + O<br />

π<br />

ln √ T<br />

, (1.3)<br />

where A is a certain constant depending on the level N and κ is the number <strong>of</strong><br />

cusps <strong>of</strong> Γ0(N).<br />

However it is widely believed (by the Sarnak-Phillips philosophy, see for<br />

example [80]) that for a “generic” non-cocompact but c<strong>of</strong>inite Fuchsian group<br />

Γ, unless there are some arithmetic or geometric symmetries present, there<br />

should only be at most a finite discrete spectra (generic here is in the sense that<br />

in some appropriate deformation space the exceptional groups have measure<br />

0).<br />

The <strong>Maass</strong> <strong>waveforms</strong> are also more mysterious than the holomorphic automorphic<br />

functions since there are only very few examples <strong>of</strong> explicit formulas<br />

for constructing <strong>Maass</strong> <strong>waveforms</strong>, whereas there are numerous examples <strong>of</strong><br />

explicit formulas for holomorphic modular functions.<br />

19


1.2 Some Structural Theory <strong>of</strong> M(Γ0(N), χ,λ)<br />

For the rest <strong>of</strong> this section we put Γ = Γ0(N) for brevity and we also assume<br />

that χ is a Dirichlet character. To facilitate the computation <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong><br />

we need as much information as possible about the various symmetries<br />

that can be used. First we will see that in the case <strong>of</strong> a real character we can<br />

assume that our functions are real-valued. Then we will consider the obvious<br />

translational symmetry which makes the Fourier expansion possible, and<br />

the reflectional symmetry which simplifies the said Fourier expansion. Then<br />

we will describe the less obvious symmetries, the Hecke operators and the<br />

involutions, which further refine the spectral eigenspaces.<br />

1.2.1 The Conjugation Operator<br />

Let f ∈ M(Γ, χ,λ) and consider the conjugation operator K, K f = f . It is<br />

clear that K commutes with the Laplacian so K f is also an eigenfunction <strong>of</strong> ∆<br />

with the same eigenvalue as f . We also see that for A ∈ Γ we get<br />

(K f ) |A(z) = K f (Az) = f (Az)<br />

= χ(A) f (z)<br />

= χ(A)(K f )(z),<br />

so K f is automorphic with respect to the conjugate <strong>of</strong> χ, i.e. K f ∈ M(Γ, χ,λ).<br />

In particular, if χ is a real character, then ¯f = K f ∈ M(Γ, χ,λ). By considering<br />

the functions 1<br />

1<br />

2 ( f + K f ) and 2i ( f − K f ) (which are real-valued) we get<br />

the following proposition.<br />

Proposition 1.2.1. If χ is a real-valued character then M(Γ, χ,λ) has a C−linear<br />

basis consisting <strong>of</strong> real-valued functions.<br />

1.2.2 The Fourier Series<br />

Notation 1.2.1. In connection with Fourier series <strong>of</strong> functions in M(Γ, χ,λ)<br />

we will use the following notation:<br />

and<br />

e(x) = e 2πix ,<br />

κn(y) = κn(R,y) = √ yKiR(2π|n|y),<br />

where λ = 1<br />

4 + R2 . Observe that κ−n(y) = κn(y).<br />

If f ∈ M(Γ, χ,λ), we define functions f j = f |σ j and we know that f j(z+1) =<br />

f (σ jSz) = f (S jσ jz) = χ(S j) f j(z) = f j(z) if the character leaves the cusp p j<br />

open.<br />

20


Let λ > 0. By separating variables in the Laplace equation it follows that f j<br />

has a Fourier expansion <strong>of</strong> the form ∑n=0 c j(n)κn(y)e(nx) (cf. [39], and recall<br />

that f is necessarily a cusp form, cf. Section 1.1.4).<br />

Proposition 1.2.2. If χ is regular for Γ then a function f ∈ M(Γ, χ,λ) admits<br />

a Fourier expansion at each cusp <strong>of</strong> Γ, and these expansions are given by the<br />

functions f j. Explicitly we have the following expansion at the cusp p j :<br />

f j = f |σ j = ∑ c j(n)κn(y)e(nx), j = 1,...,κ.<br />

n=0<br />

1.2.3 Involutions and Normalizers<br />

Definition 1.2.1. For Γ ⊆ PSL(2,R) we say that g ∈ PGL(2,R) is a normalizer<br />

<strong>of</strong> Γ in PGL(2,R) if<br />

gΓg −1 = Γ.<br />

This means in particular that A ∈ Γ ⇒ ∃B ∈ Γ such that gA = Bg. But we<br />

need a stronger definition if a character is present.<br />

We say that g is a normalizer <strong>of</strong> (Γ, χ) if g is a normalizer <strong>of</strong> Γ and for all<br />

A ∈ Γ:<br />

χ(gAg −1 ) = χ(A).<br />

The set <strong>of</strong> all normalizers <strong>of</strong> Γ in PGL(2,R) forms a group, the so-called<br />

normalizer subgroup <strong>of</strong> Γ in PGL(2,R).<br />

In the above definition we have PGL(2,R) = GL(2,R)/C, where C is the<br />

center in GL(2,R) and consists <strong>of</strong> all diagonal matrices <strong>of</strong> the form <br />

a 0<br />

0 a with<br />

a = 0.<br />

We define an action <strong>of</strong> GL(2,R) on H as follows, for g = <br />

a b<br />

c d and z ∈ H:<br />

⎧<br />

az+b<br />

⎪⎨ cz+d , if ad − bc > 0,<br />

g(z) =<br />

⎪⎩ az+b<br />

cz+d , if ad − bc < 0.<br />

Note that all <strong>of</strong> C acts trivially on H, and hence we also obtain a welldefined<br />

action <strong>of</strong> PGL(2,R) on H. And if g is a normalizer <strong>of</strong> (Γ, χ) and<br />

f ∈ M(Γ, χ,λ) then f |g ∈ M(Γ, χ,λ).<br />

The usual definition <strong>of</strong> a (linear) involution <strong>of</strong> a vector space S is a linear<br />

operator T : S → S such that T 2 = Id (or equivalently T −1 = T ).<br />

Definition 1.2.2. A linear operator T : M(Γ, χ,λ) → M(Γ, χ,λ) will be called<br />

a (Γ, χ)-involution if T 2 is the identity operator, that is<br />

for any function f ∈ M(Γ, χ).<br />

T 2 f = f<br />

21


Example 1.2.1. There are three typical examples <strong>of</strong> (Γ, χ)−involutions that<br />

we will use.<br />

1. All W ∈ PSL(2,R) which are normalizers <strong>of</strong> (Γ, χ) and satisfy W 2 ∈ Γ and<br />

χ(W 2 ) = 1 are (Γ, χ)-involutions when acting via f ↦→ f |W . In particular<br />

we will use ωN : z ↦→ −1<br />

Nz , which can be represented by<br />

0 −1<br />

√N<br />

√ N 0<br />

PSL(2,R). <br />

2. If J =<br />

1<br />

0<br />

0<br />

−1<br />

is a normalizer <strong>of</strong> (Γ, χ) then f ↦→ f |J is a (Γ, χ)involution.<br />

Note that this is just a reflection in the imaginary axis, J(z) =<br />

−¯z.<br />

It is important not to confuse the cusp normalizing map defined in the beginning<br />

with the normalizer <strong>of</strong> (Γ, χ) defined above. But some cusp normalizing<br />

maps might be normalizers <strong>of</strong> (Γ, χ) and even involutions <strong>of</strong> M(Γ, χ). In general,<br />

if it is possible, this is a very good choice <strong>of</strong> cusp normalizing map as we<br />

will see.<br />

1.2.4 The Reflection Operator<br />

By looking at the fundamental domain <strong>of</strong> PSL(2,Z) (or the Ford fundamental<br />

domain <strong>of</strong> any congruence subgroup) we see that there is an obvious symmetry;<br />

reflection in the imaginary axis. This symmetry provides us with a<br />

partitioning <strong>of</strong> the spectrum into even and odd functions and amounts actually<br />

to looking separately at Dirichlet or Neumann boundary conditions when we<br />

view the spectral problem on the fundamental domain itself (which has finite<br />

boundary points) and not on the corresponding Riemann surface PSL(2,Z)\H<br />

(with the only boundary point at the cusp).<br />

As remarked in Example 1.2.1, the reflection operator is represented<br />

by the<br />

1 0<br />

a b<br />

matrix J =<br />

, and for any matrix T = we define<br />

0 −1<br />

c d<br />

T ∗ = JT J −1 <br />

1 0 1 0 a −b<br />

=<br />

T<br />

=<br />

. (1.4)<br />

0 −1 0 −1 −c d<br />

Then J is a (Γ, χ)−involution whenever T ∗ ∈ Γ and χ(T ∗ ) = χ(T ) for all T ∈<br />

Γ. In particular this is clearly true for Γ0(N) when χ is a Dirichlet character<br />

since then χ(T ∗ ) = χ(d) = χ(T ).<br />

Suppose now that J is a (Γ, χ)−involution. We can diagonalize M(Γ, χ,λ)<br />

with respect to J and the eigenvalues will be either 1 or −1. f ∈ M(Γ, χ,λ) is<br />

said to be even or odd, respectively, if f |J = f or f |J = − f . Every even or odd<br />

22<br />

<br />


function f in M(Γ, χ,λ) has a cosine resp. sine Fourier series. This can be<br />

seen as follows. Let<br />

f (z) = ∑ c(n)κn(y)e(nx)<br />

n=0<br />

with κn(y) as in Notation 1.2.1 on page 20 (remember that κ−n(y) = κn(y)).<br />

Then<br />

f |J(z) = ∑ c(n)κn(y)e(−nx)<br />

n=0<br />

= ∑ c(−n)κn(y)e(nx).<br />

n=0<br />

If f |J = f (−¯z) = ε f , then we get c(−n) = εc(n) so<br />

f (z) = ∑ c(n)κn(y)e(nx)<br />

n=0<br />

=<br />

=<br />

=<br />

∞<br />

∑<br />

n=1<br />

κn(y)(c(n)e(nx) + c(−n)e(−nx))<br />

∞<br />

∑ c(n)κn(y)(e(nx) + εe(−nx))<br />

n=1<br />

∞<br />

∑ c(n)κn(y)<br />

n=1<br />

<br />

2cos(2πnx), ε = 1,<br />

2isin(2πnx), ε = −1.<br />

Set a(n) = 2c(n) and b(n) = 2ic(n), then f will have a Fourier cosine or sine<br />

series with coefficients a(n) or b(n) respectively, i.e.<br />

f (z) =<br />

f (z) =<br />

∞<br />

∑<br />

n=1<br />

a(n)κn(y)cos(2πnx), or<br />

∞<br />

∑ b(n)κn(y)sin(2πnx).<br />

n=1<br />

So we know that f = f0 can always be taken as an eigenfunction <strong>of</strong> J, but we<br />

should also note that this need not be the case with the other Fourier series, f j.<br />

In the next section we give a condition for the simultaneous diagonalization<br />

<strong>of</strong> all f j with respect to J.<br />

1.2.5 Complete Symmetrization<br />

Since it is usually preferable to work with symmetrized Fourier series containing<br />

only cosines or sines instead <strong>of</strong> the exponential functions we make a<br />

comment on when this is possible.<br />

23


Theorem 1.2. Let u<br />

v be any cusp satisfying (u,v) = 1, v|N and v, N<br />

<br />

v = 1 or<br />

2. Then if we take p j = u<br />

v as a cusp representative, there is a choice <strong>of</strong> cusp<br />

normalizing map σ j such that if f ∈ M(Γ, χ,λ) is an eigenfunction <strong>of</strong> J with<br />

eigenvalue ε then f j is also an eigenfunction <strong>of</strong> J with the eigenvalue χ(d)ε,<br />

where d is unique modulo N determined by d ≡ 1 mod v and d ≡ −1 mod N<br />

v .<br />

Pro<strong>of</strong>. (We refer to [51, p. 36] for the basic facts).<br />

We know that every cusp for Γ has a representative <strong>of</strong> the form u<br />

v where v|N<br />

and (u,v) = 1, and u u′<br />

v and v ′ are equivalent cusps if and only if v = v ′ and u ≡<br />

u ′ mod v, N<br />

<br />

u<br />

u<br />

v . Hence v is equivalent to − v if and only if 2u ≡ 0 mod v, N<br />

<br />

v<br />

and this is possible if and only if v, N<br />

<br />

v = 1 or 2. Suppose now that this is<br />

the case, and that we have p j = u<br />

v for such choice <strong>of</strong> u, and that A = <br />

a b<br />

Nc d ∈<br />

Γ0(N) is such that Ap j = −p j. One verifies by a straightforward computation<br />

that up to a sign change (A ←→ −A), the most general form <strong>of</strong> A is given<br />

by taking c as an arbitrary integer satisfying <br />

N<br />

v cu ≡ −2 mod v (thus c is<br />

uniquely determined modulo v if v, N<br />

<br />

v = 1, and uniquely determined modulo<br />

v<br />

2 if v, N<br />

<br />

v = 2), and then letting<br />

d = a = −1 − N<br />

u(1 − d)<br />

· cu, b = .<br />

v v<br />

In particular, note that d ≡ −1 mod <br />

N and d ≡ 1 mod v.<br />

v<br />

Observe now that f |T |J = f |J|T ∗ (where T ∗ is as in (1.4)). This means in<br />

particular that if f has J−eigenvalue ε, then f j|J = f |σ j|J = f |J|σ ∗ = ε f<br />

j |σ ∗. We<br />

j<br />

can choose the cusp normalizing map <strong>of</strong> p j as<br />

√m<br />

u x<br />

0<br />

σ j =<br />

v y 0 √ m −1<br />

<br />

, (1.5)<br />

for some integers x,y satisfying uy − vx = 1 and m = N<br />

(N,v2 is the width <strong>of</strong> the<br />

)<br />

cusp p j (note that m = N N<br />

v or 2v if v, N<br />

<br />

v = 1 or 2). This choice <strong>of</strong> σ j is unique<br />

up to right multiplication with Sr for some r ∈ R. Then we will have<br />

σ ∗ j = Jσ jJ −1 √m<br />

u −x<br />

0<br />

=<br />

−v y 0 √ m −1<br />

<br />

,<br />

which is clearly a cusp normalizing map for −p j. Since Ap j = −p j, then<br />

another cusp normalizing map <strong>of</strong> −p j is the map Aσ j, and hence we have<br />

σ ∗ j = Aσ jSt for some t ∈ R.<br />

It follows that f j|J = ε f |σ ∗ = ε f<br />

j |Aσ jSt = χ(d)ε f j|St , and hence if t ∈ Z we<br />

have f j|J = χ(d)ε f j, as desired.<br />

Hence it only remains to make sure that we can obtain t ∈ Z above. In general<br />

this can always be achieved by replacing σ j by ˜σ j where σ˜ j = σ jS t 2 , since<br />

24


then ˜σ ∗ j = σ ∗ j S− t 2 = Aσ jSt S− t 2 = Aσ jS t 2 = A ˜σ j. However, if either v, N<br />

<br />

v = 1<br />

or m is odd, then we can take x and y to be integers in (1.5). To see this, note<br />

that by writing out the upper right element in the matrix relation σ ∗ j = Aσ jSt (and using au + bv = u) one obtains −x = umt + ax + by. Using the formulas<br />

for a and b this implies that mt = N d−1<br />

v cx+ v y. Hence t ∈ Z holds if and only if<br />

N d − 1<br />

cx + y ≡ 0 mod m. (1.6)<br />

v v<br />

Recall that m = N N<br />

1−d<br />

v or 2v , so that (1.6) is equivalent to v y ≡ 0 mod m. Hence<br />

we achieve our goal if we let x,y be any integers satisfying uy − vx = 1 and<br />

y ≡ 0 (mod m). This is clearly possible if either v, N<br />

<br />

v = 1 or m is odd, since<br />

then (v,m) = 1. (In the remaining case, i.e. v, N<br />

<br />

N<br />

v = 2 and m = 2v is even, one<br />

can check that one must allow x or y to be non-integral to obtain t ∈ Z.)<br />

Corollary 1. Suppose that N can be written as N = 2r p1 ··· pm, with 0 ≤ r ≤ 3,<br />

and if m > 0 then p1,..., pm are distinct odd primes, and suppose also that<br />

in every character decomposition χ = χvχ N , v|N, the character χ N is even.<br />

v<br />

v<br />

Then we can take all f j as simultaneous eigenfunctions <strong>of</strong> J with the same<br />

eigenvalue.<br />

Note that if we do not have simultaneous eigenfunctions <strong>of</strong> J we have to<br />

work (numerically) with exponentials instead <strong>of</strong> sines and cosines, and complex<br />

instead <strong>of</strong> real Fourier coefficients.<br />

1.2.6 Hecke Operators<br />

It is well-known that the classical theory <strong>of</strong> Hecke operators (cf., e.g. [5], [36]<br />

etc.) can be carried over from the case <strong>of</strong> holomorphic modular forms to the<br />

case <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong>. This section provides an outline <strong>of</strong> the theory for<br />

those not familiar with Hecke theory, and it will also serve as a recollection <strong>of</strong><br />

fundamental facts and a common ground <strong>of</strong> notation for the more experienced.<br />

The Hecke operators considered here are operators acting on spaces <strong>of</strong> modular<br />

functions (i.e. functions automorphic with respect to congruence subgroups).<br />

We will define the Hecke operators in a similar way as in Atkin-Lehner [5].<br />

Let N ∈ Z + be given. For a prime p there is a subgroup <strong>of</strong> finite index in<br />

Γ0(N),<br />

<br />

<br />

a b<br />

<br />

Γ0(N, p) =<br />

∈ Γ0(N) <br />

c d<br />

b ≡ 0 mod p .<br />

For any d ∈ Z + , we define the map Ad : z ↦→ dz in PSL(2,R). For d = p<br />

we then have A −1<br />

p Γ0(N, p)Ap = Γ0(N p) ⊆ Γ0(N) and it is easy to see that<br />

25


χ(A −1<br />

p BAp) = χ(B) so that Ap gives a map between the spaces <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong><br />

A −1<br />

p : M(Γ0(N), χ) → M(Γ0(N, p), χ), (1.7)<br />

f ↦→ fp = f |A −1<br />

p .<br />

We then want to map fp to a function once again in M(Γ0(N), χ) but in a<br />

non-trivial way (i.e. not via Ap). Let µ<br />

R j be the set <strong>of</strong> right coset repre-<br />

j=1<br />

sentatives <strong>of</strong> Γ0(N, p) in Γ0(N) used in [5, lemma 5]. R j will then have the<br />

lower right entry d ≡ 1 mod N. If V is any element in Γ0(N), then for each<br />

j there exists a unique i such that R jV = gRi for some g ∈ Γ0(N, p), and different<br />

j’ s give different i’s. Using our special choice <strong>of</strong> R j one also checks<br />

that χ(V ) = χ(g) in this relation. Thus for any function h ∈ M(Γ0(N, p), χ)<br />

it is clear that the function ∑ µ<br />

j=1 h |R j belongs to M(Γ0(N), χ). Hence we may<br />

define the operator Tp : M(Γ0(N), χ) → M(Γ0(N), χ) as follows:<br />

Tp f = 1<br />

√ p<br />

µ<br />

∑ fp|R j<br />

j=1<br />

.<br />

The factor <strong>of</strong> 1<br />

√p is a convenient normalization. Working out the coset representatives<br />

explicitly (see [5, lemma 5]) gives the following formula for a<br />

prime number p with (N, p) = 1<br />

Tp f (z) = 1<br />

p−1<br />

√<br />

p ∑ f<br />

j=0<br />

z + j<br />

p<br />

Tq f (z) = Uq f (z) = 1 √ q<br />

<br />

+ 1<br />

√ p χ(p) f (pz),<br />

and for prime q with q|N (compare, e.g., [69, p. 142 (4.5.26)]):<br />

<br />

z + j<br />

.<br />

q−1<br />

∑ f<br />

j=0<br />

We will follow the convention to call Tq with q|N an exceptional Hecke operator<br />

and denote it by Uq instead <strong>of</strong> Tq. In principle, due to the multiplicative<br />

nature <strong>of</strong> the Hecke operators (cf. Theorem 1.3 below) it is sufficient to define<br />

Hecke operators for primes, and then use the multiplicativity (cf. 1.9 below<br />

and [69, lemma 4.5.12]) to recursively define Tn for any positive integer n<br />

with (n,N) = 1. To use the above construction for a composite number n is<br />

quite elaborate and it is not done this way in the literature. Hecke (cf. [36, nos.<br />

32, 35, 36, 37 and 41, p. 859]) does not explain how he arrives at the exact<br />

definition <strong>of</strong> the operators Tn. To get a motivation for the formula (1.8) below<br />

(including the character) it is easiest to use the approach <strong>of</strong> double cosets as<br />

in Shimura ([89]) or Miyake ([69]). Using either <strong>of</strong> these constructions yields<br />

the following definition.<br />

26<br />

q


Definition 1.2.3. For any n ∈ Z + and any f ∈ M(Γ0(N), χ,λ) we define the<br />

Hecke operator Tn by the formula<br />

Tn f (z) = 1 <br />

az + j<br />

√<br />

n ∑<br />

. (1.8)<br />

d<br />

ad=n<br />

d>0<br />

d−1<br />

χ(a) ∑ f<br />

j=0<br />

Theorem 1.3. The Hecke operators Tn with (n,N) = 1 are endomorphisms <strong>of</strong><br />

the space M(Γ0(N), χ,λ) which have the following properties:<br />

TnTm = ∑<br />

d|(m,n)<br />

d>0<br />

χ(d)Tmn , (1.9)<br />

d2 for any integers n and m with (n,N) = (m,N) = 1, and (for the adjoint operator)<br />

T ∗<br />

n = χ(n)Tn. (1.10)<br />

In particular <br />

Tn<br />

(n,N) = 1 forms a commutative family <strong>of</strong> normal operators<br />

which also commutes with the Laplacian and the reflection J.<br />

Pro<strong>of</strong>. To prove the multiplicative property cf. [89, eq. (3.5.10)] and look at<br />

the action <strong>of</strong> Tn on the Fourier coefficients <strong>of</strong> a <strong>Maass</strong> waveform f . Suppose<br />

that c(k) and b(k) are the respective Fourier coefficients <strong>of</strong> f and Tn f , then the<br />

following relation holds:<br />

b(k) = ∑<br />

ad=n,a,d>0<br />

χ(a)c<br />

<br />

kd<br />

, (1.11)<br />

a<br />

with the usual convention that c(r) = 0 if r /∈ Z. We prove the equality (1.9)<br />

by comparing the action on the Fourier coefficients <strong>of</strong> the left hand side and<br />

the right hand side.<br />

Next, to prove (1.10), a quick computation using (1.9) shows that it suffices<br />

to treat the case n = p a prime. To prove that T ∗ p = χ(p)Tp one uses the<br />

following relation between the Peterson inner products on the two groups:<br />

< Tp f ,g > Γ0(N)= 1<br />

√ p < fp,g > Γ0(N,p) .<br />

This relation is an easy extension <strong>of</strong> [5, lemma 12] to the case <strong>of</strong> non trivial<br />

character using the fact that χ(R j) = 1 for all right coset representatives R j <strong>of</strong><br />

Γ0(N, p) in Γ0(N) as given in [5, lemma 5]. Alternatively: see [51, thm. 6.20]<br />

and [69, thm. 4.5.4(1)]. It is easy to show that Tn commutes with ∆ and J.<br />

Remark 1.2.1. Actually ([69, thm. 4.5.13]) we also have<br />

UqTn = TnUq,<br />

27


for q|N and (n,N) = 1 but in general the operators Uq are not normal so they<br />

are put aside for a while, until the next section, where we will introduce the<br />

space <strong>of</strong> newforms.<br />

Theorem 1.4. There exists an orthogonal basis {φ j} in M(Γ0(N), χ,λ) consisting<br />

<strong>of</strong> eigenfunctions to all Hecke operators Tn with (n,N) = 1.<br />

Pro<strong>of</strong>. Observe that M(Γ0(N), χ,λ) is a finite dimensional vector space (see<br />

[39, pp. 140, 298]) so we can choose an ON basis { f j} m j=1 and represent the<br />

operators Tn by matrices Λ(n) in this basis. From Theorem 1.3 above it is<br />

clear that Λ(n) <br />

(n,N) = 1 is a family <strong>of</strong> commuting normal m × m matrices.<br />

From elementary linear algebra arguments (the spectral theorem or [84,<br />

thm. 9.4.1]) we know that there exists a basis {φ j} m j=1 where each φ j is an<br />

eigenfunction <strong>of</strong> all Λ(n), and this is exactly what we want. Compare: [51,<br />

thm. 6.21] and [69, thm. 4.5.4(3)].<br />

1.2.7 Oldforms and Newforms<br />

From (1.11) it is clear that if f has Fourier coefficients c(n) and Tn f = λ(n) f ,<br />

then<br />

c(n) = λ(n)c(1).<br />

We would like to make sure that c(1) = 1 here, but unfortunately as we will<br />

see below, it may happen that c(1) = 0.<br />

Suppose that χ is an imprimitive character mod N <strong>of</strong> conductor q. If q|M<br />

and Md|N for some positive integers M,d, let χM be the character induced<br />

from χ on Γ0(M). Then<br />

f (z) ∈ M(Γ0(M), χM,λ) ⇒ f (z), f |Ad (z) ∈ M(Γ0(N), χ,λ)<br />

(with Ad as in section 1.2.6, cf. (1.7)), which is clear from the fact that<br />

AdΓ0(N)A −1<br />

<br />

N<br />

d = Γ0 d ,d <br />

N ⊆ Γ0 d ⊆ Γ0(M), and χ(AdBA −1<br />

d ) = χ(B).<br />

Definition 1.2.4. Let f ∈ M(Γ0(N), χ,λ). We say that f is an oldform if there<br />

exist positive integers M and d with Md|N, a character χM mod M and a<br />

function f1 ∈ M(Γ0(M), χM,λ) such that we can write f = f 1|A d . The smallest<br />

such integer M is called the conductor (or level) <strong>of</strong> f . Compare: [51, p. 107<br />

(bottom)] and [69, p. 162 (bottom)].<br />

28<br />

Suppose now that f is an oldform with Fourier series given by<br />

f (z) = ∑ c(m)<br />

m=0<br />

√ yKiR(2π|m|y)e(mx).


Then<br />

f |Ad (z) = f (dz) = ∑ c(m)<br />

m=0<br />

dyKiR(2π|m|dy)e(mdx)<br />

= √ d ∑<br />

n≡0 mod d<br />

c<br />

<br />

n<br />

<br />

√yKiR(2π|n|y)e(nx).<br />

d<br />

Hence if d > 1 then the first Fourier coefficient <strong>of</strong> f (dz) is zero. This is an<br />

example <strong>of</strong> one <strong>of</strong> the inconveniences <strong>of</strong> oldforms.<br />

Let M old (Γ0(N), χ,λ) be the linear space spanned by the oldforms and then<br />

define M new (Γ0(N), χ,λ) as the orthogonal complement with respect to the<br />

Petersson inner product. We then have<br />

M(Γ0(N), χ,λ) = M old (Γ0(N), χ,λ) ⊕ M new (Γ0(N), χ,λ).<br />

It is easy to see that both subspaces are stable under the Hecke operators Tn<br />

with (n,N) = 1, since T ∗<br />

n = χ(n)Tn.<br />

Definition 1.2.5. A function f ∈ M new (Γ0(N), χ,λ) is called a normalized<br />

newform <strong>of</strong> level N if f is an eigenfunction <strong>of</strong> J and all Hecke-operators Tn,<br />

(n,N) = 1, and the first Fourier coefficient <strong>of</strong> f is 1, i.e. c(1) = 1.<br />

Proposition 1.2.3. If χ is a primitive character mod N and q|N, then the operator<br />

Uq is unitary on M(Γ0(N), χ,λ); i.e.<br />

U ∗ q = U −1<br />

q .<br />

Pro<strong>of</strong>. The pro<strong>of</strong> is elementary (at least for prime indices N) and goes by first<br />

proving that U ∗ q = ωNUqω −1<br />

N and then using properties <strong>of</strong> the Petersson innerproduct.<br />

Here ωN is defined as in Example 1.2.1. (For a sketch <strong>of</strong> the pro<strong>of</strong><br />

see [12]).<br />

The following theorem is crucial to the theory <strong>of</strong> oldforms/newforms. See<br />

Atkin-Lehner [5] and Miyake [69] for similar results in the classical setting.<br />

For a more general adelic version, see Miyake [68]. The following version for<br />

<strong>Maass</strong> <strong>waveforms</strong> is in principle a direct adaptation <strong>of</strong> [94, thm. 4.6] to the<br />

case <strong>of</strong> non-trivial character.<br />

Theorem 1.5. Let χ be a Dirichlet character with conductor q. The spaces <strong>of</strong><br />

newforms have the following properties.<br />

(a) For each multiple N <strong>of</strong> q, the normalized newforms <strong>of</strong> level N and<br />

character χ, mN Fj , form an orthogonal basis in the newspace<br />

j=1<br />

Mnew (Γ0(N), χ,λ). If we want to stress the level and/or character<br />

we write Fj = F (N,χ)<br />

j . This is called the Hecke-(eigen)basis <strong>of</strong><br />

Mnew (Γ0(N), χ,λ).<br />

29


(b) Each Fj is an eigenfunction <strong>of</strong> all Tn with n ∈ Z + .<br />

(c) A basis in the space M(Γ0(N), χ,λ) is given by<br />

<br />

F (M,χ)<br />

<br />

<br />

i (dz) <br />

d,M ∈ Z+ <br />

,q|M ,dM|N, 1 ≤ i ≤ mM .<br />

(d) Let χ ′ be a Dirichlet character mod M. If f ∈ M new (Γ0(N), χ,λ)<br />

and g ∈ M new (Γ0(M), χ ′ , µ) with λ,µ > 0, then either there exists<br />

an infinite number <strong>of</strong> primes p for which the Hecke eigenvalues <strong>of</strong><br />

Tp are different for f and g, or else N = M, λ = µ, χ ′ = χ and<br />

f ≡ g.<br />

Pro<strong>of</strong>. Observe that the only thing that differs from [94, thm. 4.6] is that we<br />

have a non-trivial character.<br />

The pro<strong>of</strong> <strong>of</strong> (a)-(c) can be done exactly as in the pro<strong>of</strong> <strong>of</strong> [94, thm. 4.6]. The<br />

pro<strong>of</strong> <strong>of</strong> “multiplicity one”, (d) is a bit more complicated to extend. The main<br />

difference from the case <strong>of</strong> trivial character is the fact that ωN (cf. Example<br />

1.2.1) and Tn no longer commute and we need the Euler products <strong>of</strong> both f<br />

and ωN f for the pro<strong>of</strong> <strong>of</strong> Lemma 4.12 in [94], however it is a simple matter<br />

to do the same thing as in the pro<strong>of</strong> <strong>of</strong> [69, thm. 4.6.19], noting that on the<br />

newspace both Tp and Uq are normal so any eigenform <strong>of</strong> all Hecke operators<br />

Tn is also an eigenform <strong>of</strong> all T ∗<br />

n and U ∗ q (i.e T ∗<br />

n = ω −1<br />

N TnωN and U ∗ q = U −1<br />

q on<br />

the newspace).<br />

Remark 1.2.2. Note that Uq is not normal on the full space M(Γ0(N), χ,λ)<br />

(unless, <strong>of</strong> course, if χ is primitive, in which case there are only newforms) so<br />

therefore we would have no multiplicity one theorem if we tried to diagonalize<br />

the full space with respect to all Hecke operators.<br />

Theorem 1.6. Let f be a normalized newform with Tn f = λ(n) f and f |J = ε f .<br />

If the Fourier expansion <strong>of</strong> f is<br />

(where c(1) = 1) then<br />

f (z) = ∑<br />

|m|≥1<br />

c(m) √ yKiR(2πmy)e(mx)<br />

c(m) = λ(m), and<br />

c(−m) = ελ(m),<br />

for all m ∈ Z + . We also have the following multiplicativity relations<br />

30<br />

c(m)c(n) = ∑<br />

d|(m,n),d>0<br />

χ(d)c<br />

mn<br />

d 2<br />

c(m)c(p) = c(mp), p|N, m ∈ Z.<br />

<br />

, (n,N) = 1, m ∈ Z, (1.12)


Pro<strong>of</strong>. The first part follows by adding f and f |J and reordering the sum. The<br />

second part follows at once from the multiplicative relation in Theorem 1.3<br />

which then apply to the Hecke eigenvalues λ(n), and when we set c(1) = 1 in<br />

the relation c(n) = λ(n)c(1) we are done.<br />

Proposition 1.2.4. Suppose that f ∈ M new (Γ0(N), χ,λ) is a normalized newform,<br />

and that q is a prime such that q|N, then the following holds.<br />

(a) If χ is primitive then |c(q)| = 1.<br />

(b) If χ is trivial and q 2 ∤ N then c(q) = λq<br />

√q , where λq = ±1. In fact<br />

λq is the eigenvalue <strong>of</strong> −Wq, cf. Subsection 1.2.8 below.<br />

(c) If χ is trivial and q 2 |N, then c(q) = 0.<br />

Pro<strong>of</strong>. For (a), use Prop. 1.2.3. The assertions (b) and (c) are the analogs <strong>of</strong><br />

[5, thm. 3]. Compare: [69, thm. 4.6.17 and and cor. 4.6.18(2)], and [84, thm.<br />

9.4.8].<br />

Recall that by Theorem 1.5 we can always choose a Hecke eigenbasis in the<br />

newspace M new (Γ0(N), χ,λ) where the functions are normalized by c(1) = 1.<br />

We stress that if the character χ is non-trivial, then we can not in general<br />

assume that the Hecke basis is real valued. For by Theorem 1.3 we have<br />

T ∗<br />

n = χ(n)Tn whenever (n,N) = 1, and thus<br />

λ(n) = χ(n)λ(n) (1.13)<br />

and c(n) = χ(n)c(n). Hence, in particular, if χ(n) = −1 then c(n) is purely<br />

imaginary, whereas if χ(n) = 1 then c(n) is real.<br />

Remark 1.2.3. Suppose that χ is real and that f is a normalized newform in<br />

the space M new (Γ0(N), χ,λ) with f |J = ε f and Hecke eigenvalues λ(n). Then<br />

ε ¯f is also a normalized newform in M new (Γ0(N), χ,λ), but with Hecke eigenvalues<br />

λ(n). Now f and ε ¯f are in general linearly independent since λ(n) =<br />

λ(n) whenever χ(n) = 1 and λ(n) = 0. Thus the newspace M new (Γ0(N), χ,λ)<br />

is in general multidimensional.<br />

An exception to the above fact are the CM-type forms à la [46] (originally<br />

constructed by <strong>Maass</strong> in [62]). These forms are real and (hence) they have<br />

c(n) = 0 for all n ∈ Z + with (n,N) = 1 and χ(n) = 1.<br />

31


1.2.8 The Cusp Normalizing Maps as Normalizers <strong>of</strong> (Γ0(N), χ)<br />

<br />

Suppose that Q|N and = 1. Then following [5] and [6] we can attach<br />

Q, N<br />

Q<br />

a normalizer <strong>of</strong> Γ0(N) to Q, call it WQ . This can be taken <strong>of</strong> the form<br />

<br />

WQ =<br />

Qx<br />

Nz<br />

y<br />

Qw<br />

,<br />

where x,y,z,w ∈ Z, x ≡ 1 mod N<br />

Q , y ≡ 1 mod Q and Q2xw − Nzy = Q. One<br />

can also verify that WQ is a normalizer <strong>of</strong> (Γ0(N), χ) for certain χ (for examples<br />

see below). A remarkable fact about WQ is that as an operator on<br />

M(Γ0(N), χ,λ) it is independent <strong>of</strong> the choices <strong>of</strong> x,y,z and w. In fact, if<br />

W ′ Q =<br />

<br />

Qx ′ y ′ <br />

for any x ′ ,y ′ ,z ′ ,w ′ ∈ Z and detW ′ Q = Q, then (as operators)<br />

Nz ′ Qw ′<br />

W ′ Q = χQ(y ′ )χ N Q<br />

any f ∈ M(Γ0(N), χ,λ)<br />

(x ′ )WQ (see [6, prop. 1.1]). In particular, this implies that for<br />

f |J|WQ = χQ(−1) f |WQ|J.<br />

For us the normalizer WQ will be most useful for real characters because <strong>of</strong><br />

the following proposition (cf. [6, prop. 1.1]).<br />

<br />

Proposition 1.2.5. Suppose that Q|N and Q, N<br />

<br />

Q = 1, and let χ = χQχ N .<br />

Q<br />

Then for any f ∈ M(Γ0(N), χ,λ) we have WQ f = f |WQ ∈ M(Γ0(N), χQχ N ,λ),<br />

Q<br />

and W 2 Q f = χQ(−1)χ N (Q) f . In particular, if χQ is real and χQ(−1)χ N (Q) = 1<br />

Q<br />

Q<br />

then WQ maps M(Γ0(N), χ,λ) into itself and WQ is a (Γ, χ)−involution.<br />

The normalizer WQ can in many cases be taken to be a cusp normalizing<br />

map, as explained in the following proposition.<br />

<br />

Proposition 1.2.6. Suppose that Q|N and = 1. Then we can choose<br />

WQ as a cusp normalizing map <strong>of</strong> the cusp Q<br />

N .<br />

Q, N<br />

Q<br />

In reference to Definition 1.1.4, note that WQ is equal to Q − 1 2WQ ∈ PSL(2,R).<br />

Pro<strong>of</strong>. We may choose x = z = 1 in the formula for WQ, that is<br />

<br />

WQ =<br />

Q<br />

N<br />

y<br />

wQ<br />

,<br />

where y,w ∈ Z, and Qw − N<br />

Q<br />

Qy = 1. Note that the width <strong>of</strong> the cusp N is<br />

32<br />

<br />

N,<br />

N<br />

N<br />

Q<br />

2 =<br />

N<br />

Q<br />

N<br />

<br />

Q, N<br />

Q<br />

= Q


can be written<br />

uniquely up to a translation as σ = Aρ, where A ∈ SL(2,Z) maps ∞ to Q<br />

N , and<br />

ρ is a scaling by the width <strong>of</strong> the cusp. Hence, writing WQ as<br />

<br />

1 y Q 0<br />

WQ =<br />

,<br />

0 1<br />

(cf. [51, p. 37]). We know that the cusp normalizing map <strong>of</strong> Q<br />

N<br />

N<br />

Q wQ<br />

we see that WQ is a cusp normalizing map for Q<br />

N .<br />

Another interesting property <strong>of</strong> WQ is the following (cf. [6, prop. 1.2, 1.3]).<br />

Proposition 1.2.7. Let Q ∈ Z + <br />

be such that Q|N and<br />

and p be primes with (p,N) = 1, q|N and (q,Q) = 1. Then<br />

TpWQ f = χQ(p)WQTp f , and<br />

UqWQ f = χQ(q)WQUq f ,<br />

Q, N<br />

Q<br />

<br />

= 1, and let q<br />

for f ∈ M(Γ0(N), χ,λ), where χQ is the character mod Q induced by χ (i.e. χ =<br />

χQχ N ).<br />

Q<br />

Observe that for a Hecke eigenfunction f we get Tp(WQ f ) = χQ(p)WQTp f =<br />

χQ(p)λpWQ f where λp is the Hecke eigenvalue <strong>of</strong> f for Tp. In other words<br />

WQ f is also an eigenfunction <strong>of</strong> Tp with eigenvalue λpχQ(p). In view <strong>of</strong> “mul-<br />

tiplicity one” (cf. Thm. 1.5 (d)), if χQ ≡ 1 and f is a newform, then WQ f must<br />

be a multiple <strong>of</strong> f , i.e. WQ f = µQ f , where µ 2 Q = χ N (Q) (cf. Prop. 1.2.5).<br />

Q<br />

Now if we choose the cusp normalizer σ j <strong>of</strong> Q<br />

N as σ j = WQ (cf. Prop. 1.2.6),<br />

then we have f j = µ j f1 , and in terms <strong>of</strong> the different Fourier expansions, this<br />

means that (for all k)<br />

c j(k) = µ jc1(k), (1.14)<br />

where µ 2 j = χ N (Q). This relation can either be used as a means to test the<br />

Q<br />

accuracy <strong>of</strong> our program or to reduce the size <strong>of</strong> the linear system in Section<br />

1.3.2 (cf. the normalization (B) on p. 38). However, it is important to observe<br />

that (1.14) does not in general allow us to reduce the Fourier coefficients at all<br />

cusps to the ones at i∞. Cf. Prop. 1.2.6.<br />

Our comment in Prop. 1.2.4(b) about −Wq is justified by combining the<br />

foregoing remarks with the three references cited in the pro<strong>of</strong> <strong>of</strong> Prop. 1.2.4.<br />

Cf. [6, §1,2], [51, pp. 113–118] and [69, cor. 4.6.18] for some related ideas.<br />

33


1.3 <strong>Computational</strong> <strong>Aspects</strong><br />

1.3.1 Introduction<br />

We know that a <strong>Maass</strong> waveform f ∈ M(Γ0(N), χ,λ), with λ > 0 is completely<br />

described by its Fourier series at ∞, but unfortunately to assure stability<br />

<strong>of</strong> the numerical method we need knowledge <strong>of</strong> the Fourier series at all<br />

cusps <strong>of</strong> F, i.e. for 1 ≤ j ≤ κ,<br />

f j(z) = ∑ c j(n)κn(y)e(nx).<br />

|n|≥1<br />

The aim here is to compute such functions f , that is, we wish to find an eigenvalue<br />

λ = 1<br />

4 +R2 (for convenience we will usually refer to R as the eigenvalue)<br />

and a set <strong>of</strong> Fourier coefficients {c j(n)|1 ≤ j ≤ κ, |n| ≥ 1}. In general (the<br />

exceptions are the CM-type forms, cf., e.g., [46]) there are no known formulas<br />

for neither the eigenvalues nor the coefficients and we have to be content with<br />

numerical approximations.<br />

There are two steps (Phase 1 and Phase 2) in the algorithm. First we locate R<br />

(up to the desired accuracy) and an initial set <strong>of</strong> Fourier coefficients, and then,<br />

if we want to, we can use this initial set to generate a larger set <strong>of</strong> coefficients.<br />

1.3.2 Phase 1<br />

The method here is a generalization <strong>of</strong> Hejhal’s algorithm in [43] to groups<br />

with several cusps and non-trivial characters, and some <strong>of</strong> the ideas are also<br />

inspired by the algorithm by Selander and Strömbergsson in [95]. The idea is<br />

that given a real number R we use linear algebra to compute a set <strong>of</strong> numbers<br />

that are likely to be close to “true” Fourier coefficients if R is close to a “true”<br />

eigenvalue <strong>of</strong> a <strong>Maass</strong> waveform.<br />

Preliminary Numerical Remarks<br />

First we introduce an “effective zero”. In standard double precision arithmetic<br />

we know that x + 10 −16 x ≈ x, so if we fix ε < 10 −16 then a truncation error <strong>of</strong><br />

ε can in principle be neglected (computationally). We will use [[ε]] to denote<br />

a quantity with absolute value less than ε.<br />

Suppose that f ∈ M(Γ0(N), χ,λ) with Fourier expansions at the cusps<br />

f j(z) = ∑ c j(n)κn(y)e(nx), 1 ≤ j ≤ κ. (1.15)<br />

|n|≥1<br />

There is a trivial (a priori) bound on the coefficients: c j(n) = O( √ n) (cf. [39,<br />

p. 585] and [50, thm. 3.2]). Combining this with κn(y) = √ yKiR (2π|n|y) and<br />

the fact that KiR(u) ∼ π<br />

2ue−u , as u → ∞ (for fixed R), we see that the tail<br />

<strong>of</strong> the sum (1.15) satisfies ∑|n|≥M = O e−2πyM as M → ∞, with the implied<br />

34


constant depending on R and y. Hence for any fixed y (and R) we can take<br />

M = M(y) such that<br />

f j(z) = ∑<br />

1≤|n|≤M(y)<br />

c j(n)κn(y)e(nx) + [[ε]], (1.16)<br />

for 1 ≤ j ≤ κ. Using similar observations as in [43, p. 311] it is clear that, for<br />

R large, we can take M(y) as<br />

<br />

R + AR<br />

M(y) =<br />

1 <br />

3<br />

, (1.17)<br />

2πy<br />

for some constant A. In practice it turns out that A = 12 or 15 is good enough.<br />

We will now see how to use the automorphy condition ((1.19) below) to get a<br />

good linear system that can be solved for the c j(n).<br />

The Linear System<br />

In order to make the algorithm stable we need to use the Fourier coefficients<br />

at all cusps, which means that we have to use different expansions at different<br />

regions in a clever way as follows (see also [95]). Recall the notation from<br />

Subsection 1.1.2.<br />

We will view f ∈ M(Γ0(N), χ,λ) as pieced together <strong>of</strong> the Fourier series<br />

at all cusps, meaning that we use the f j that is most convenient at each point.<br />

Hence, given w ∈ FN we will use the identity<br />

f (w) = χ(U −1<br />

w ) f I(w)(σ −1<br />

I(w) Uww). (1.18)<br />

Cf. Remark 1.1.1 on p. 15.<br />

Given any point z ∈ H and any j ∈ {1,...κ}, we let z j = σ jz (which may<br />

or may not be inside FN), and let w j be the pullback <strong>of</strong> z j into FN, i.e. w j =<br />

Tj(z j) ∈ FN with Tj ∈ Γ0(N), and let z∗ −1<br />

j = σI( j) Uw jw j (here I( j) = I(w j)), then<br />

z∗ −1<br />

j = σI( j) Uw jTjσ jz, so<br />

f j(z) = f (σ jz) = f (T −1 −1<br />

j Uw j σI( j)z ∗ j)<br />

= χ(T −1 −1<br />

j Uw j ) fI( j)(z ∗ j). (1.19)<br />

This relation (implicit automorphy, cf. [43, p. 298 (15)]) between f j(z) and<br />

fI( j)(z∗ j ) is what enables us to get hold <strong>of</strong> all Fourier series. Consider the<br />

truncated Fourier series<br />

ˆf j(z) = ∑<br />

1≤|n|≤M(y)<br />

c j(n)κn(y)e(nx).<br />

One way to view this series is as a Discrete Fourier Transform, and we can<br />

perform an inverse transform over the following set <strong>of</strong> sampling points along<br />

35


a horocycle:<br />

<br />

<br />

<br />

zm = xm + iY <br />

xm = 1<br />

<br />

1<br />

(m − ), 1 − Q ≤ m ≤ Q ,<br />

2Q 2<br />

for some Y < Y0 (recall (1.2) on page 15), and Q > M(Y ). The inverse transform<br />

gives us that for 1 ≤ |n| ≤ M(Y ) < Q we have<br />

c j(n)κn(Y ) = 1<br />

2Q<br />

= 1<br />

2Q<br />

Q<br />

∑<br />

m=1−Q<br />

Q<br />

∑<br />

m=1−Q<br />

ˆf j(zm)e(−nxm)<br />

f j(zm)e(−nxm) + [[ε]],<br />

where in the last step we used (1.16). This system is now almost a tautology,<br />

but we can use the implicit automorphy (1.19) to get a good “mix” <strong>of</strong> the<br />

Fourier coefficients and a far from trivial system.<br />

Mimicking the the discussion leading to (1.18) we let Tm j ∈ Γ0(N) be the<br />

pullback map <strong>of</strong> σ jzm into FN, and let wm j = Tm j(σ jzm) ∈ FN, I(m, j) =<br />

−1<br />

= σI(m, j) Um jwm j. Furthermore we let χm j =<br />

). Using (the analog <strong>of</strong>) (1.19) we then obtain:<br />

I(wm j), Um j = Uwm j ∈ Γ0(N), z∗ m j<br />

χ(T −1 −1<br />

m j Um j<br />

c j(n)κn(Y ) = 1<br />

2Q<br />

Q<br />

∑ χm j fI( j,m)(z<br />

m=1−Q<br />

∗ m j)e(−nxm) + [[ε]],<br />

and we now want to substitute the truncated Fourier expansion for fI( j,m).<br />

Since we know that ℑ(z∗ −1<br />

m j ) = ℑ(σI( j,m) Um jTm jσ j(zm j)) ≥ Y0 we can actually<br />

use the same truncation point, M0 = M(Y0) for all series, and in interchanging<br />

the order <strong>of</strong> summation we get the following expression, valid for 1 ≤ |n| ≤<br />

M(Y ) < Q and 1 ≤ j ≤ κ :<br />

where<br />

c j(n)κn(Y ) =<br />

V ji 1<br />

nk =<br />

2Q<br />

Q<br />

∑<br />

m=1−Q<br />

I( j,m)=i<br />

κ<br />

∑<br />

∑<br />

i=1 1≤|k|≤M0<br />

ci(k)V ji<br />

nk + 2[[ε]], (1.20)<br />

χm jκk(y ∗ m j)e(kx ∗ m j)e(−nxm). (1.21)<br />

Remark 1.3.1. It seems clear that in order for the system (1.20) to be useful,<br />

we need to use the automorphy relation non-trivially sufficiently <strong>of</strong>ten. In<br />

practice we found that a necessary (and sufficient!) condition for the numerics<br />

to behave well is that z ∗ m j = zm for all j,m. This condition is ensured if we<br />

36


keep ℑ(z∗ m j ) > Y for all j,m. This is <strong>of</strong> course automatically fulfilled if we<br />

choose Y < Y0 as above, but in many cases it is possible to verify numerically<br />

that the same inequality also holds for certain choices <strong>of</strong> Y (slightly) larger<br />

than Y0. This means that we may be able to use a Y > Y0 and a corresponding<br />

M(Y ) < M0. The drawback is that such a Y can not be used for Phase 2 later<br />

since there the pullbacks z∗ m j <strong>of</strong> the horocyclic points zm will eventually fill out<br />

the whole fundamental domain (as Q → ∞ and Y → 0 simultaneously).<br />

We now have a linear system that can be used to obtain the coefficients.<br />

Note that the V ji<br />

nk can be small due to either the K-Bessel decay or “bad mixing”<br />

in the sense that the numbers <strong>of</strong> m such that I( j,m) = i might differ<br />

much (meaning that the pullback <strong>of</strong> the horocycle does not encircle each cusp<br />

equally much). Bad mixing is avoided basically through increasing the length<br />

<strong>of</strong> the horocycle by means <strong>of</strong> decreasing Y . The system (1.20) can be expressed<br />

as<br />

0 =<br />

κ<br />

∑<br />

∑<br />

i=1 |k|≤MY<br />

ci(k) ˜V ji<br />

nk + 2[[ε]], (1.22)<br />

where ˜V ji ji<br />

nk = Vnk − δnkδ jic jκn(Y ). The term −κn(Y ) occurring in all diagonal<br />

entries gives us good reason to hope that this system should turn out to be<br />

well-conditioned, just as in [43, p. 298 (19) and second paragraph <strong>of</strong> p. 299]).<br />

We now have a linear system which is satisfied by (linear combinations <strong>of</strong>)<br />

Fourier coefficients <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong> f ∈ M(Γ0(N), χ,λ) (up to the error<br />

2[[ε]]). If we let V denote the (κM0 × κM0)−matrix ˜V ji<br />

nk and C denote the<br />

κM0−vector <strong>of</strong> Fourier coefficients c j(k), we can write the system as<br />

VC = 0. (1.23)<br />

Observe that the solution space <strong>of</strong> this system for a true eigenvalue R is at least<br />

a 1-dimensional linear space, so in order to get a unique solution we need to<br />

use some sort <strong>of</strong> normalization.<br />

Of course, if f and all the various f j are all Fourier cosine series, there is<br />

an immediate “cosine” counterpart <strong>of</strong> (1.23). Similarly for sine series. This<br />

remark is important in all cases where Cor. 1 applies. The matrix coefficients<br />

in these analogs <strong>of</strong> (1.23) are all real!<br />

Normalization<br />

Depending on the dimension <strong>of</strong> the particular space M(Γ0(N), χ,λ) we are<br />

investigating, we have to employ different types <strong>of</strong> normalizations algorithmically.<br />

In order to test the accuracy <strong>of</strong> our coefficients (cf. p. 41 below) we<br />

will always try to get eigenfunctions <strong>of</strong> Hecke operators and/or involutions<br />

whenever possible. We will now discuss two examples (A) and (B) <strong>of</strong> possible<br />

normalizations.<br />

37


(A) Newforms, trivial character. According to a standard conjecture (and<br />

experimental findings) the space <strong>of</strong> newforms with trivial character should<br />

have dimension 1. So, by Theorem 1.5 we can search for a normalized newform<br />

by setting c1(1) = 1 and dropping the first equation from V . If we find a<br />

true R-value under this assumption, the coefficients will automatically satisfy<br />

all multiplicative relations (cf. Theorem 1.6).<br />

(B) N and real primitive χ satisfying the hypotheses <strong>of</strong> Cor. 1. To discuss<br />

this case some preliminaries are necessary. We present them in outline form.<br />

1) Let MR(Γ0(N), χ,λ) ⊆ M(Γ0(N), χ,λ) be the space <strong>of</strong> all real-valued<br />

<strong>Maass</strong> <strong>waveforms</strong>. By simple use <strong>of</strong> the operator K, one sees that<br />

M(Γ0(N), χ,λ) = MR(Γ0(N), χ,λ) + iMR(Γ0(N), χ,λ)<br />

as vector spaces. The operators WQ and Tn are all endomorphisms on the<br />

spaces M(Γ0(N), χ,λ) and MR(Γ0(N), χ,λ) (see Prop. 1.2.5 and Thm. 1.3).<br />

One knows that W 2 Q = ±Id, and that TnWQ = χQ(n)WQTn for (n,N) = 1 (can<br />

be proved by induction on Prop. 1.2.7.)<br />

3) When χ is primitive, there are no oldforms. Cf. Thm. 1.5. Thm. 1.5(d)<br />

assures us <strong>of</strong> multiplicity one for any newforms.<br />

4) Apply Thm. 1.4 to M(Γ0(N), χ,λ) and get an orthogonal basis F1,...,Fm<br />

consisting <strong>of</strong> normalized newforms. (Do not confuse things here with (1.15).)<br />

By Thm. 1.6 and Cor. 1 each Fj is either a cosine or sine series at infinity and<br />

remains <strong>of</strong> the same type at the other cusps.<br />

5) It is now natural to apply Remark 1.2.3 to the basis functions Fj. If we<br />

form the the normalized newform εFj there are two possibilities.<br />

a) We arrive back at Fj. This is equivalent to saying λ(n) = λ(n), for all n.<br />

Since λ(n) = χ(n)λ(n), for (n,N) = 1, by (1.13), this would imply λ(n) = 0<br />

whenever χ(n) = −1.<br />

b) We arrive at a form independent <strong>of</strong> Fj. By mult. 1, we thus have that<br />

the sequences λ(n) and λ(n) are essentially different as n → ∞. Accordingly,<br />

here, via (1.13), we have λ(n) =nonzero pure imaginary for infinitely many<br />

primes with χ(n) = −1. While, λ(1) = 1, <strong>of</strong> course.<br />

6) We apply 5) to each Fj, and thus deduce a natural splitting <strong>of</strong> the basis<br />

{F1,...,Fm} into singlets and pairs.<br />

7) In pursuing the numerics on the machine, we hypothesize (at least initially)<br />

that either m = 1 (singlet) or m = 2 (one pair) for our given space<br />

M(Γ0(N), χ,λ). In the following comments, we assume ε = 1. The treatment<br />

<strong>of</strong> ε = −1 is similar.<br />

8) Take m = 1 first. Look at 1<br />

2F1. This is a unit normalized real cosine<br />

series. Its expansions at the other cusps are again real cosine series. Take the<br />

cosine analog <strong>of</strong> (1.23) and declare λ(1) = 1. Solve the system in the usual<br />

way with two parallel Y. This case should have a unique solution.<br />

38


Cases where χ(n) = −1 should yield coefficients that are exactly 0. This<br />

will be very compelling to see! (Cf. Table 1.5, right column.)<br />

9) Now assume m = 2. The (complex) basis is {F1,F2}, where F2 = εF1 =<br />

F1 is the “Remark 1.2.3 flip” <strong>of</strong> F1. The idea here is to construct a real basis<br />

from F1 and F2. To this end we construct F+ = 1<br />

4 (F1 + F2) and F− =<br />

1<br />

4i (F1 − F2). These functions are now real cosine series at all cusps. F+ is<br />

unit normalized, while the n = 1 term for F− is zero. Observe that the Fourier<br />

coefficients <strong>of</strong> F+ and F− are ℜ(λ(n)) and ℑ(λ(n)) respectively (here λ(n)<br />

are the coefficients <strong>of</strong> F1).<br />

The “cosine” analog <strong>of</strong> (1.23) holds. To single out F+, we impose the condition<br />

that λ(n) = 0 for χ(n) = −1. We reduce the size <strong>of</strong> the “cosine” system<br />

(1.23) accordingly, and set λ(1) = 1, and use our two parallel Y. This process<br />

determines λ (the eigenvalue) and F+. If we instead impose the condition<br />

λ(n) = 0 for χ(n) = 1, and set λ(p0) = 1 where χ(p0) = −1 and use the appropriately<br />

reduced “cosine” system with parallel Y , we obtain aF− with some<br />

nonzero a ∈ R. To get F1 and F2, we form 2F+ ± 2i <br />

1<br />

a aF− and determine a<br />

either by Prop. 1.2.4(a) or by use <strong>of</strong> relation (1.12). If desired, λ can again be<br />

solved for at this stage <strong>of</strong> the game.<br />

10) An alternative reduction in the size <strong>of</strong> the “cosine” system (1.23) can be<br />

achieved by looking at the action <strong>of</strong> WQ on M(Γ0(N), χ,λ). Cf. Prop. 1.2.7<br />

and the discussion afterwards. (This discussion can be generalized to (n,N) =<br />

1 by induction; notice too that there are no oldforms when χ is primitive.)<br />

Each WQ (Fj) must be a nonzero multiple <strong>of</strong> some Fk by virtue <strong>of</strong> its Hecke<br />

action. Depending on whether m = 1 or 2, the resulting permutation structure<br />

must either be (1), (12), or (1)(2). In each instance, the relation W 2 Q = ±Id<br />

affords one some a priori control on any “non zero constants” that arise.<br />

Let {WQ1 ,...,WQr } be any commutative family <strong>of</strong> WQ’s satisfying W 2 Q = Id<br />

and the hypotheses <strong>of</strong> Prop. 1.2.5. Cf. [6, prop. 1.4]. Since each WQ acts<br />

unitarily, it is natural to seek a basis <strong>of</strong> MR(Γ0(N), χ,λ) consisting <strong>of</strong> simultaneous<br />

eigenfunctions <strong>of</strong> the WQ j . By Prop. 1.2.6, the Fourier expansion <strong>of</strong><br />

any such basis function at the cusp Q j<br />

N will be “redundant”; there is thus a<br />

corresponding reduction in the size <strong>of</strong> the “cosine” analog <strong>of</strong> (1.23).<br />

Remark 1.3.2. If Cor. 1 applies, the reduction with respect to the WQ’s, as<br />

mentioned above, can be applied also to the case <strong>of</strong> the trivial character in a<br />

similar manner.<br />

Remark 1.3.3. Note that in the case <strong>of</strong> N = 4 and trivial character, using<br />

Thm. 1.2 and the fact that the cusp normalizing map <strong>of</strong> the cusp p3 = − 1<br />

2 ,<br />

σ3 : z ↦→ z<br />

2z+1 is in the normalizer <strong>of</strong> Γ0(4) we can actually use pure sine/cosine<br />

series and real coefficients at all cusps in this case too. And we can also use<br />

reduction with respect to σ2 = ωN and σ3. This is not true in general for nonsquare-free<br />

levels, e.g. Γ0(9). Then we must use complex arithmetic.<br />

39


Example 1.3.1. As an illustration <strong>of</strong> (B) consider the case <strong>of</strong> prime level N ≡<br />

1 mod 4 and χ = <br />

N . With the notation as in (B), for the sake <strong>of</strong> simplicity,<br />

·<br />

suppose that m = 2 and ε1 = ε2 = 1. Let F1 and F2 = F1 be the Hecke eigenbasis<br />

<strong>of</strong> normalized newforms in M(Γ0(N), χ,λ). There is only one WQ, i.e. ωN,<br />

which is now also a cusp normalizing map, and ω2 N = Id. We will see how to<br />

diagonalize with respect to ωN to get a reduction <strong>of</strong> the linear system. Observe<br />

that Tn (ωNF1) = χ(n)ωNTnF1 = χ(n)λ(n)ωNF1, and hence by multiplicity one<br />

we must either have ωNF1 = ±F1 (since ω2 N = Id) or we have ωNF1 = µF2 and<br />

ωNF2 = µ −1F1, for some constant µ ∈ S1 (again since ω2 N = Id). In case<br />

that ωNF1 = ±F1 we already have an eigenfunction <strong>of</strong> ωN in F1 and we don’t<br />

have to do anything more. If µ = ±1 then the functions F±(defined under (B))<br />

are also invariant under ωN, i.e. ωNF± = ±F± if µ = 1 and ωNF± = ∓F± if<br />

µ = −1. If µ = ±1, then we can form<br />

f± = 1<br />

1 ± µ (F1 ± µF2).<br />

It is easily verified that f± are both unit normalized (have first Fourier coefficient<br />

equal to 1), real-valued on H, and satisfy ωN f± = ± f±. Denote<br />

the Fourier coefficients <strong>of</strong> f+ and f− by c + (n) and c − (n), respectively. For<br />

(n,N) = 1, one immediately checks that<br />

and<br />

c + (n) = c − (n) = λ(n) for χ(n) = +1<br />

c + (n)c − (n) = λ(n) 2 ≤ 0 for χ(n) = −1.<br />

To find f+, we use the “cosine” analog <strong>of</strong> (1.23). We reduce the size <strong>of</strong> the<br />

system by using the fact that ωN is a cusp normalizing map, i.e. we set<br />

c2(k) = c1(k), k = 1,...,M0<br />

together with c1(1) = 1. And then we drop the corresponding equations. Similarly<br />

for f−. In this case, when we run with the parallel Y, we have a lot <strong>of</strong><br />

extra tests; Hecke relations coming from the λ(n), c + (n) = c − (n) if χ(n) = 1<br />

and c + (n)c − (n) ≤ 0 if χ(n) = −1. All these can be used as independent tests<br />

for the accuracy <strong>of</strong> our program.<br />

By running through the search for λ both for f+ and f− we get further<br />

insurance.<br />

Locating the Eigenvalues<br />

We will see how, with the use <strong>of</strong> the normalized linear system introduced<br />

above, we can locate an eigenvalue R. Cf. [44, p. 5]<br />

Consider the linear system (1.23). The entries <strong>of</strong> the coefficient matrix<br />

40


V = V (R,Y ) are clearly real-analytic functions <strong>of</strong> R, but as functions <strong>of</strong> Y<br />

they might be only piecewise continuous. For a value <strong>of</strong> R corresponding to<br />

a <strong>Maass</strong> waveform f we expect the solution vector C = C(R,Y ) to give us<br />

the actual Fourier coefficients if the appropriate normalization is used (e.g.<br />

c(1) = 1), and Y is kept less than Y0 for the sake <strong>of</strong> well-conditioning.<br />

This fact is used to determine whether a given R is close to a true eigenvalue<br />

or not by solving for C(R,Y ) and looking for known properties <strong>of</strong> Fourier<br />

coefficients <strong>of</strong> a <strong>Maass</strong> waveform. The three most important properties (here)<br />

are:<br />

a) Invariance under change <strong>of</strong> Y . As described above, for a true<br />

eigenvalue the solution vector will be invariant when we change<br />

Y (keeping Y < Y0).<br />

b) Hecke relations. We can look for Hecke eigenfunctions by seeking<br />

solutions which satisfy multiplicative relations (cf. (1.12)),<br />

for example the relation c1(2)c1(3) = c1(6).<br />

c) Involutions <strong>of</strong> (Γ0(N), χ). We use cusp normalizing maps which<br />

are involutions with eigenvalues, µ j = ±1, and seek solution vectors<br />

which satisfy c j(k) = µ jc1(k) (cf. (1.14), Example 1.3.1, and<br />

B item 10).<br />

For a given interval I = [R1,R2], we want to find all <strong>Maass</strong> <strong>waveforms</strong> with<br />

eigenvalues R in this interval. The idea here is to do this by solving the linear<br />

system (1.23) for two different values <strong>of</strong> Y in parallel and then try to find R ∈ I<br />

such that the solution vectors C(R,Yi) match (usually we only require the first<br />

few coefficients to match, and use the rest only as an insurance).<br />

In practice we first divide I into a number <strong>of</strong> equally sized small chunks:<br />

R1 = x0 < ··· < x j < ··· < xN1 = R2,<br />

where the number N1 is chosen in a way that it is probable that any interval<br />

Ij = [x j,xj+1] contains at most one eigenvalue, using e.g. Weyl’s law (1.3) and<br />

some assumption on the behaviour <strong>of</strong> the nearest-neighbour spacings (e.g. that<br />

they are exponentially distributed).<br />

We then look at each small Ij to see if there is a change <strong>of</strong> sign in any<br />

<strong>of</strong> the differences ck − c ′ k , k = 2,3,4, where ck = ck(R,Y1) and c ′ k = c′ k (R,Y2)<br />

are entries in C(R,Y1) and C(R,Y2) respectively. If there are sign changes<br />

for all considered differences in an interval Ij0 we go to the next stage <strong>of</strong> the<br />

investigation and “zoom in” into this interval. At this point we usually form a<br />

functional<br />

h(R) = ω2(c2 − c ′ 2) + ω3(c3 − c ′ 3) + ω4(c4 − c ′ 4),<br />

41


and try to minimize it over this interval. We use ω j ∈ {±1} to “align” the<br />

differences so that h(R) changes sign where all three differences change sign.<br />

The minimization can be done in a number <strong>of</strong> ways. One very efficient approach<br />

is to use the method <strong>of</strong> false position to get successively better approximations<br />

to the location <strong>of</strong> the minimum (which if it exists is near a point<br />

where h changes sign).<br />

When a value <strong>of</strong> R that approximates a zero <strong>of</strong> h is found it is listed as a<br />

candidate for a true eigenvalue. These candidates are stored and subjected to<br />

further examination.<br />

Remark 1.3.4. Zeros <strong>of</strong> the K-Bessel functions in the left hand side <strong>of</strong> (1.20)<br />

(or “random noise”) can trigger false indications <strong>of</strong> zeros in the intervals Ij0 .<br />

Keeping track <strong>of</strong> how fast c j − c ′ j changes over each interval, most <strong>of</strong> these<br />

“false” intervals are immediately discarded.<br />

If we only use two Y -values and a few coefficients, it can not be excluded that<br />

the located minimum <strong>of</strong> h is not close to a real eigenvalue, but such mistakes<br />

will be spotted either when we try to refine the eigenvalue or when we look<br />

closer at properties a)-c) for a larger set <strong>of</strong> coefficients.<br />

Eventually we declare that the R which we have found is close to an eigenvalue<br />

<strong>of</strong> the Laplacian, and that our {c j(n)} are close to the Fourier coefficients<br />

<strong>of</strong> the corresponding <strong>Maass</strong> waveform. It is worth stressing that, strictly<br />

speaking, these assertions are never proved rigorously through our computations<br />

— although the excellent agreement seen when testing properties a) -<br />

c) on a large number 1 <strong>of</strong> coefficients on top <strong>of</strong> those used in the functional<br />

h(R) clearly gives a very strong heuristic justification. The question <strong>of</strong> giving<br />

rigorous pro<strong>of</strong>s is dealt with in a forthcoming paper [15].<br />

Remark 1.3.5. To speed up the process <strong>of</strong> looking for sign changes over the<br />

large number <strong>of</strong> small intervals Ij we use Lagrange interpolation to evaluate<br />

the matrix V at this point (usually interpolation <strong>of</strong> degree 14 is good enough).<br />

Cf. [44, p. 6].<br />

It should be remarked that the algorithm reproduces both the known oldforms<br />

from Γ0(1) (see [45] or [42]), the newforms on Γ0(5) as in [9], as well<br />

as the explicitly known CM-forms (cf. [46]).<br />

There is obviously no guarantee that we find all eigenvalues in a specified<br />

interval in this manner, but comparing the results with the detailed version <strong>of</strong><br />

Weyl’s law (1.3) might give us an indication <strong>of</strong> missing eigenvalues (cf. [9,<br />

§A.1] and [8, §8]).<br />

1 e.g. 50-10000 produced by the Phase 2 algorithm, cf. Section 1.3.3<br />

42


The Pullback Algorithm<br />

Since the implicit automorphy (1.19) plays an important part in the algorithm,<br />

it is crucial to have an efficient means to compute the pullback z ∗ ∈ FN <strong>of</strong> a<br />

point z ∈ FN (observe that the notation z ∗ differs from the one used earlier).<br />

We recall that in the case <strong>of</strong> the modular group Γ0(1), it is easy to make a<br />

pullback to the standard fundamental domain, F1 = {z = x + iy ∈ H||x| ≤<br />

1<br />

1<br />

2 , |z| ≥ 1} using a sequence <strong>of</strong> “flip-flops” through the generators E : z ↦→ − z<br />

and S : z ↦→ z + 1 (cf. [84, p. 51] or [9, pp. 44-46]).<br />

Instead <strong>of</strong> extending this algorithm to the case <strong>of</strong> Γ0(N) by using side pairing<br />

generators <strong>of</strong> a suitable fundamental domain, we will use the facts that<br />

Γ0(N) is a subgroup <strong>of</strong> finite index in PSL(2,Z), and that it is easy to find a<br />

set <strong>of</strong> coset representatives.<br />

Let {Vj} vN<br />

j=1 and FN = ∪Vj(F1) be as on p. 14. Given z ∈ FN we make a<br />

pullback into F1; ˜z = T (z) ∈ F1 with T ∈ Γ0(1). Then we find the index j such<br />

that T −1 ∈ Γ0(N)Vj, and note that VjT ∈ Γ0(N). Hence the Γ0(N) pullback <strong>of</strong><br />

z is given by<br />

z ∗ = VjT (z) ∈ FN.<br />

This gives a pullback algorithm for any N, but we need the coset representatives<br />

{Vj} vN<br />

vN<br />

j=1 to apply it. The observation that {Vj} j=1 is by definition a<br />

<br />

maximal set <strong>of</strong> vN = N ∏p|N 1 + p−1 maps in PSL(2,Z), all independent<br />

over Γ0(N) allows us to use the following simple recursive algorithm:<br />

"Traverse Γ0(1) as a tree in S, E and S −1 and collect maps independent over<br />

Γ0(N) until exactly vN independent maps have been found."<br />

This calculation is done once and for all for each group and the resulting representatives<br />

are stored. A possible alternative here is to use explicit formulas<br />

for the coset representatives, but the advantage <strong>of</strong> the above algorithm is that<br />

it can be extended to other congruence subgroups, such as Γ1(N), Γ 0 (N) or<br />

Γ(N) as well. Actually it can be extended to any finite-index subgroup <strong>of</strong><br />

PSL(2,Z) that possesses some precise characteristic which we can use to see<br />

if a given map belongs to the group or not.<br />

1.3.3 Phase 2<br />

After Phase 1 is completed we have an (approximate) eigenvalue R and a<br />

corresponding set <strong>of</strong> (approximate) Fourier coefficients<br />

<br />

c j(k) <br />

1 ≤ j ≤ κ, 1 ≤ |k| ≤ M0 .<br />

43


Suppose now that we want to compute c j(n) for NA ≤ n ≤ NB. Going back to<br />

the identity (1.20),<br />

c j(n)κn(Y ) =<br />

κ<br />

∑<br />

∑<br />

i=1 |k|≤M0<br />

ci(k)V ji<br />

nk + 2[[ε]]<br />

(valid for Y < Y0, 1 ≤ |n| ≤ M(Y ) < Q and 1 ≤ j ≤ κ), we see that in order<br />

to compute c j(n) for n ∈ [NA,NB] (where NB > M0), we take Y < Y0 such that<br />

M(Y ) > NB. To avoid cancellation we also have to make sure that the particular<br />

choice <strong>of</strong> Y does not make κn(Y ) too small. We can now formulate the core<br />

<strong>of</strong> Phase 2 as a theorem.<br />

Theorem 1.7. Suppose that f ∈ M(Γ0(N), χ, 1<br />

4 + R2 ) has Fourier coefficients<br />

c j(n) . Let Y < Y0, and let n,Q be integers such that Q > M(Y ) and 1 ≤<br />

|n| ≤ M(Y ). We then have<br />

c j(n) = ∑κ ji<br />

i=1 ∑|k|≤M0 ci(k)Vnk + Errn(Y ), (1.24)<br />

κn(Y )<br />

where V ji<br />

nk is given by (1.21) on p. 36, and the error Errn(Y ) is given by<br />

Errn(Y ) = [[ε]]<br />

κn(Y ) .<br />

In order to use Theorem 1.7 to successfully compute the Fourier coefficients<br />

we have to adjust the value <strong>of</strong> Y (and accordingly M(Y ) and Q also) to keep<br />

the error Errn(Y ) from growing.<br />

If we just want to confirm the existence <strong>of</strong> an eigenvalue in the neighbourhood<br />

<strong>of</strong> R, or if we want to make a picture <strong>of</strong> an eigenfunction, the first set <strong>of</strong><br />

coefficients from Phase 1 is usually sufficient. But, even for those purposes,<br />

the method <strong>of</strong> Phase 2 is a very cheap way to improve the last part <strong>of</strong> the<br />

already obtained coefficients, as well as obtaining a lot more.<br />

1.3.4 Remarks on the Performance <strong>of</strong> the Algorithm<br />

The main difference (with regard to performance issues) between groups with<br />

one cusp, e.g. the modular group or Hecke triangle groups, and the groups<br />

Γ0(N) with N > 1 should now be clear. The presence <strong>of</strong> extra cusps introduces<br />

more sets <strong>of</strong> coefficients, thus increasing the size <strong>of</strong> the linear system used in<br />

Phase 1. We will give some examples <strong>of</strong> how the relevant factors scale with<br />

respect to N (and R).<br />

44


Table 1.1: Time to compute compute the coefficient vector for R = 9.5336... on a<br />

3.2GHz CPU<br />

Timing for a Single R<br />

N M0 T (R,N)<br />

2 13 0.03s<br />

5 32 0.14s<br />

13 90 1.17s<br />

17 118 2.20s<br />

41 307 23.9s<br />

101 757 287.4s<br />

Suppose we are given a level N and a potential eigenvalue R, and that we want<br />

to compute the corresponding (minimal) set <strong>of</strong> coefficients, C = {c j(n)|1 ≤<br />

|n| ≤ M0, 1 ≤ j ≤ κ}, i.e. we have to solve the system (1.22) under some normalization<br />

assumption. Let T (R,N) be the time required to perform this task.<br />

It turns out in practice that the time required to solve the system is negligible<br />

compared to the time required to compute the coefficient matrix, and the<br />

main contribution to T (R,N) in fact comes from the K-Bessel computations.<br />

Observe that in the evaluation <strong>of</strong> V ji<br />

nk the only K-Bessel terms we need are<br />

) for 1 − Q ≤ m ≤ Q, these are independent <strong>of</strong> n and i, hence the<br />

KiR(2π|k|y ∗ m j<br />

total number <strong>of</strong> K-Bessel computations is 2Q(M0κ) ∼ 2κM2 0 (as we can take<br />

Q ∼ M0). Using Y0 ≥<br />

√ 3<br />

2N we see that, roughly M0 = [const]N(R + 1), and<br />

T (R,N) ≪ κN 2 (R + 1) 2 . (1.25)<br />

Remember that κ is dependent on the number <strong>of</strong> prime factors <strong>of</strong> N, and for<br />

a prime number κ = 2. From this formula we see that the time to compute<br />

even a single eigenfunction with a small eigenvalue increases drastically as<br />

the level N grows. No matter how much we improve the speed <strong>of</strong> computing<br />

the individual K-Bessel functions, attaining levels like N ≈ 1000 seems to<br />

be out <strong>of</strong> reach at the present time. In table 1.1 we give some examples <strong>of</strong><br />

T (R,N) for different (prime) N, with R = 9.5336... (the first eigenvalue on<br />

PSL(2,Z)).<br />

Recall steps 1-3 on p. 41. To find eigenvalues we need to compute the<br />

matrix V (R,Y ) at least 14 × 2 = 28 times for step 1 (using Lagrange interpolation<br />

<strong>of</strong> degree 14 for the function R ↦→ V (R,Y ), cf. Remark 1.3.5) and then<br />

at least 6×2 = 12 times for the step 3 (assuming that it takes approximately 6<br />

iterations with the method <strong>of</strong> false position to obtain the required precision).<br />

45


The K-Bessel Function<br />

The K-Bessel function is computed using a number <strong>of</strong> different algorithms,<br />

each with different range <strong>of</strong> applicability (cf. [42], [98] and [9]).<br />

One possible way to speed things up is to use an adaptive Lagrange interpolation<br />

method for the function x ↦→ KiR(x). However, due to the oscillatory<br />

nature <strong>of</strong> KiR(x) for small x we need a large number <strong>of</strong> interpolation points<br />

when N is large (i.e. when Y0 is small) so this type <strong>of</strong> interpolation is in general<br />

only efficient in Phase 2 or when R is large .<br />

Timing for Phase 2<br />

Since we only use the first few V ji<br />

nk with |k| ≤ M0 in Theorem 1.7 we need<br />

about 2QM0κ K-Bessel calls to compute c j(n). Remember here that for (1.24)<br />

to hold we need Q > M(Y ) > |n|, and we also need Y < const (1 + R)/n,<br />

hence, if we take Y = [const] · (1 + R)/n,then we have M(Y ) = [const] · (1 +<br />

R)/Y = [const] · n. The time to compute one coefficient c j(n) is thus roughly<br />

proportional to<br />

nκM0 = κN(R + 1)n,<br />

and hence the time to compute the first A coefficients is approximately proportional<br />

to<br />

κN(R + 1)A 2 , as A → ∞.<br />

Observe that the K-Bessel functions in V ji<br />

nk does not depend on n or i. Hence<br />

to avoid unnecessary computations we can use a vector to store all 4QκM0<br />

K-Bessel values used in (1.24). We will then only need to perform these<br />

computations when we change Y, but unfortunately we will have to deal with<br />

memory issues instead, and these can degrade the performance as much as<br />

the computations themselves. To compute a large number <strong>of</strong> coefficients we<br />

have to keep a balance between the number <strong>of</strong> computations and the size <strong>of</strong><br />

the allocated memory. In practice this is a rather difficult task.<br />

1.4 Results<br />

The focus <strong>of</strong> the current project has been to provide a robust and efficient<br />

algorithm for computing <strong>Maass</strong> <strong>waveforms</strong> on Hecke congruence subgroups<br />

Γ0(N) for any integer N. The following are our main results so far.<br />

1. Extensive lists <strong>of</strong> eigenvalues for small and prime N = 5,7,13, with either<br />

trivial<br />

<br />

character or the real Dirichlet character mod N given by the symbol<br />

N<br />

· .<br />

2. Shorter lists <strong>of</strong> eigenvalues for all N ≤ 30. More precisely, we believe our<br />

lists contain all eigenvalues with R ≤ 10 for N = 1,2,3,4,5,6,7,8,9,10, 11,<br />

46


12, 13, 14, 15, 17, 19, 21, 22, 23, 25, 26,29, R ≤ 5 for N = 16,20,24, R ≤ 3<br />

for N = 30 and R ≤ 2 for N = 18,27 and 28.<br />

3. The first eigenvalue, λ1 = 1<br />

4 + R21 , for all prime N ≤ 131.<br />

4. Large sets <strong>of</strong> Fourier coefficients for a few <strong>Maass</strong> <strong>waveforms</strong>, e.g. for N =<br />

5 the following coefficients have been computed: all c(n) with |n| ≤ 70209<br />

for R = 3.2642...; all c(n) with |n| ≤ 89263 for R = 4.8937...; all c(n)<br />

with |n| ≤ 56420 for R = 4.1032... (the first two eigenvalues correspond<br />

to the character <br />

N<br />

· and the last one corresponds to the trivial character).<br />

Work on extending these data, as well as carrying out more extensive statistical<br />

tests, is currently being pursued. In this section we will give some<br />

examples and comments on each item in the above list.<br />

Further data than that which is presented in this section is available from<br />

the author upon request.<br />

1.4.1 Eigenvalues<br />

Tables 1.6, 1.3 and 1.4 provide some examples <strong>of</strong> eigenvalues for different<br />

groups. In tables 1.3 and 1.4 the eigenvalues <strong>of</strong> the operators J and ωN are<br />

indicated in the respective columns. The eigenvalues are denoted by − for<br />

−1, + for 1, and ∗ is used to indicate that both eigenvalues are present. In all<br />

tables we use H1 = |c1(2)c1(3) − c1(6)| as one indication <strong>of</strong> the accuracy <strong>of</strong><br />

the program (we based our search on property a) on p. 41; hence no Hecke<br />

relation like H1 was built in to our algorithm).<br />

It should be remarked that in the case <strong>of</strong> Table 1.6, Γ0(4), all computed<br />

newforms are eigenfunctions with eigenvalue −1 <strong>of</strong> both non-trivial cusp normalizers,<br />

i.e. σ2 = ω4 : z ↦→ −1/4z, and σ3 : z ↦→ z/(2z + 1), which are both<br />

Γ0(4)−involutions.<br />

1.4.2 Lowest Eigenvalues<br />

We have used the Phase 1 algorithm to compute small eigenvalues for the<br />

groups M(Γ0(N), χ) as N ranges through the primes up to 131 and χ is either<br />

trivial or the real Dirichlet character <br />

N<br />

· . The first located eigenvalues are<br />

listed in Table 1.2.<br />

Recall that all eigenvalues fall into distinct classes with respect to the involutions<br />

J and ωN. We denote the classes by {++,+−,−+,−−}.<br />

In the case <strong>of</strong> a non-trivial character, CM-forms are present, and in all cases<br />

considered such a form was found to occur as the lowest eigenvalue. All the<br />

CM-forms occurring here are <strong>of</strong> the type considered in [20] (i.e. the narrow<br />

class number <strong>of</strong> Q[ √ N] is 1), and by the explicit formula (cf. [20, p. 112]) the<br />

CM-forms with lowest eigenvalues are <strong>of</strong> the type − − . In table 1.2, H1and<br />

H2 are two parameters indicating the error. H1 is always defined by H1 =<br />

47


Figure 1.1: Histogram <strong>of</strong> Fourier coefficients for <strong>Maass</strong> forms on Γ0(5), 5 ·<br />

350<br />

300<br />

250<br />

200<br />

150<br />

100<br />

50<br />

0<br />

−2 −1.5 −1 −0.5 0 0.5 1 1.5 2<br />

(a) R = 3.2642... (CM-form)<br />

250<br />

200<br />

150<br />

100<br />

50<br />

0<br />

−2 −1.5 −1 −0.5 0 0.5 1 1.5 2<br />

(b) R = 4.8937...<br />

|c1(2)c1(3) − c1(6)|, and H2 is defined as either the true error in R for CM-<br />

<br />

forms, or as H2 = |c1(N)| − 1<br />

<br />

<br />

(cf. Prop. 1.2.4) otherwise.<br />

√ N<br />

1.4.3 Fourier Coefficients<br />

Table 1.5 gives some examples <strong>of</strong> Fourier coefficients and Hecke relations<br />

obtained by the Phase 2 algorithm. Figure 1.1 (a) is a histogram <strong>of</strong> the co-<br />

efficients c1(p) <strong>of</strong> the CM-form with eigenvalue 3.26... where the first 3407<br />

primes for which χ(p) = = 1 are used. Figure 1.1 (b) is a histogram<br />

5<br />

p<br />

<strong>of</strong> the ˜c1(p) for one <strong>of</strong> the <strong>Maass</strong> <strong>waveforms</strong> on (Γ0(N), χ5) with eigen-<br />

5<br />

value R = 4.8937.... Here ˜c1(p) = c1(p) if χ(p) = p = 1 and −ic1(p)<br />

<br />

5<br />

if χ(p) = p = −1 (to get real coefficients), and the first 7462 primes are<br />

used. The experimentally obtained density functions (curves) in Figures 1.1<br />

(a) and (b) should be compared with the known density function 1<br />

<br />

π 4 − x2 − 1 2<br />

in the CM-case (a) and the conjectured density function 1<br />

√<br />

2π 4 − x2 in the non-<br />

CM case (b) (cf. [46, 45]).<br />

48


Table 1.2: Lowest eigenvalues for Γ0(N)<br />

N χ J ωN R H1 H2<br />

2 - - 5.4173348068447 0.3E-13 0.5E-12<br />

3 - - 4.38805356322 0.1E-10 0.8E-11<br />

5 - - 3.02837629306 0.4E-11 0.7E-12<br />

<br />

5· 5 - - 3.264251302636 0.7E-11 0.2E-12 CM<br />

7 - - 1.92464430511 0.2E-11 0.8E-11<br />

11 - - 2.03309099385 0.3E-12 0.1E-11<br />

13 - - 0.9708154174 0.7E-10 0.2E-10<br />

<br />

13· 13 - - 1.31473442107 0.1E-10 0.7E-12 CM<br />

17 - - 1.4414285450 0.5E-10 0.3E-10<br />

<br />

17· 17 - - 0.7498863406 0.5E-10 0.4E-10 CM<br />

19 - - 1.0919915598 0.2E-13 0.3E-10<br />

23 - + 1.3933371415 0.6E-10 0.2E-10<br />

29 + + 1.017266551 0.5E-09 0.8E-09<br />

<br />

29· 29 - - 0.953598000 0.3E-11 0.5E-10 CM<br />

31 + + 0.789356178 0.8E-10 0.4E-08<br />

37 - - 0.64230596 0.4E-08 0.6E-08<br />

<br />

37· 37 - - 0.630391294 0.7E-09 0.7E-09 CM<br />

41 + + 0.66572483 0.9E-08 0.3E-08<br />

<br />

41· 41 - - 0.37767450 0.2E-08 0.1E-07 CM<br />

43 - - 0.65545239 0.5E-09 0.7E-08<br />

47 - + 0.5854522 0.5E-07 0.1E-08<br />

53 - - 0.8039894582 0.2E-12 0.4E-08<br />

<br />

53· 53 - - 0.799094454 0.5E-10 0.8E-09 CM<br />

59 + + 0.595829688 0.4E-10 0.3E-08<br />

61 - - 0.4180624 0.4E-07 0.9E-07<br />

<br />

61· 61 - - 0.42868522 0.2E-08 0.3E-08 CM<br />

67 - - 0.67759021 0.3E-07 0.1E-08<br />

71 + + 0.35745048 0.8E-08 0.5E-08<br />

73 - - 0.517887505 0.1E-08 0.2E-09<br />

CM This is a CM-form. For definition <strong>of</strong> H2 see the end <strong>of</strong> the table<br />

49


N χ J ωN R H1 H2<br />

<br />

73·<br />

73<br />

- - 0.20488585 0.5E-07 0.7E-08 CM<br />

79 + + 0.55177485 0.2E-06 0.9E-12<br />

83 - + 0.640287578 0.3E-09 0.3E-09<br />

89 + + 0.4894360 0.1E-07 0.1E-07<br />

<br />

89· 89<br />

- - 0.2273961 0.3E-08 0.6E-07 CM<br />

97 - - 0.41693167 0.1E-07 0.8E-08<br />

<br />

97· 97<br />

- - 0.16846 0.4E-05 0.1E-05 CM<br />

101 + + 0.45375925 0.6E-09 0.3E-08<br />

103 - - 0.56540670 0.9E-08 0.5E-10<br />

107 - + 0.5816789 0.2E-07 0.8E-07<br />

109 - - 0.423655 0.2E-05 0.9E-07<br />

<br />

109 - - 0.2822870 0.8E-09 0.3E-07 CM<br />

109<br />

·<br />

113 - - 0.46200265 0.4E-09 0.3E-07<br />

<br />

113 - - 0.21379233 0.2E-08 0.8E-08 CM<br />

113<br />

·<br />

127 + + 0.373385109 0.5E-07 0.4E-10<br />

131 - + 0.261072758 0.7E-09 0.4E-11<br />

CM This is a CM-form<br />

For a CM-form H2 is the true error and for a non-CM form we know<br />

by Proposition 1.2.4(a) that |c(N)| = 1, and we put H2 = ||c(N)| − 1|.<br />

Table 1.3: Eigenvalues for N = 29, χ = <br />

29<br />

·<br />

50<br />

R J ωN H1 H2 a<br />

0.9535979998 - - 0.9E-10 0.5E-10 (CM)<br />

1.804894697229 - * 0.2E-13 0.2E-12<br />

1.907196000095 + + 0.3E-13 0.6E-12 (CM)<br />

1.9491483792 + * 0.9E-14 0.2E-10<br />

2.1682948312 - * 0.7E-13 0.3E-10<br />

2.8084567797 - * 0.8E-12 0.4e-10<br />

2.8252375274 + * 0.2E-13 0.8E-10<br />

2.8607940001 - - 0.5E-13 0.1E-11 (CM)<br />

a H2 is the true error in R for CM-forms and otherwise<br />

|c+(4) − c−(4)|. Cf. Remark 1.3.1.


Table 1.4: Eigenvalues for N = 23, χ = 1<br />

R J ωN H1 H2 = ||c(23)| − 1<br />

√ 23 | a<br />

1.3933371415 - + 0.8E-11 0.4E-11<br />

1.5061266371 - - 0.3E-11 0.1E-10<br />

1.5795892401 + + 0.7E-12 0.3E-11<br />

1.9352504926 - - 0.3E-11 0.3E-10<br />

2.4307774166 - - 0.1E-10 0.9E-10<br />

2.4643338705 + + 0.7E-11 0.9E-09<br />

2.4936237868 - + 0.2E-11 0.7E-11<br />

2.5321343295 + + 0.1E-11 0.1E-09<br />

2.6109599620 + - 0.4E-12 0.3E-11<br />

2.7362777132 - + 0.1E-10 0.1E-09<br />

2.8079149467 - - 0.6E-13 0.1E-10<br />

3.0166209534 + - 0.3E-11 0.5E-11<br />

3.2357419725 + + 0.2E-12 0.2E-10<br />

3.3723339187 - - 0.2E-10 0.8E-11<br />

3.5202562438 + - 0.1E-11 0.8E-13<br />

3.5245436383 - + 0.5E-10 0.1E-09<br />

3.5882850865 - + 0.9E-12 0.3E-11<br />

3.7062705391 - - 0.6E-10 0.2E-12<br />

3.8167552363 + - 0.3E-12 0.4E-12<br />

3.8189248762 - - 0.5E-10 0.2E-09<br />

3.8400352089 + + 0.6E-11 0.1E-11<br />

3.9122175318 + + 0.9E-10 0.7E-10<br />

4.16025275662 + + 0.5E-12 0.1E-11<br />

4.19680364676 - - 0.6E-11 0.4E-11<br />

4.23167066214 - + 0.9E-10 0.5E-09<br />

4.30661013365 - - 0.7E-11 0.8E-11<br />

4.43672242684 + + 0.4E-10 0.8E-08<br />

4.44033811150 + + 0.7E-10 0.6E-08<br />

4.50005039136 - + 0.6E-10 0.4E-11<br />

4.50100759650 + - 0.1E-11 0.2E-12<br />

4.73923930491 - + 0.1E-11 0.5E-12<br />

4.78525397008 - - 0.7E-12 0.3E-11<br />

4.83751695193 - - 0.5E-11 0.6E-11<br />

4.85889850005 + - 0.2E-10 0.8E-10<br />

4.88226240518 + - 0.3E-10 0.1E-10<br />

4.98487107442 - + 0.4E-10 0.1E-08<br />

4.99908300143 - + 0.2E-10 0.6E-10<br />

a From Proposition 1.2.4(b) we know that |c(23)| = 1<br />

√23 .<br />

51


Table 1.5: Fourier coefficients for N = 5 and χ = <br />

5<br />

·<br />

52<br />

R = 4.893781291438 R = 3.264251302636 (CM) a<br />

c(2) 1.217161411799i<br />

c(3) 0.295119713347i<br />

c(5) b exp(1.157414657530i)<br />

c(7) −1.138873755146i<br />

c(11) −0.041396292578<br />

c(13) −0.558344591841i<br />

c(17) −0.212576664102i<br />

c(19) 0.608670097807<br />

c(23) −1.205831908853i<br />

c(29) 0.162328579004<br />

c(31) −0.556019364974<br />

c(37) 0.411889174623i<br />

c(41) −0.835153489179<br />

c(43) −1.240299728133i<br />

c(47) −0.605458209042i<br />

c(53) 1.458187614411i<br />

c(59) 0.922374897976<br />

c(61) −1.247703551349<br />

c(67) 0.33772433772i<br />

c(71) −0.74219625846<br />

c(73) −0.44062515256i<br />

c(79) 0.71467631335<br />

c(83) 0.99829038428i<br />

c(1553) −1.3885884029i<br />

c(72977) −0.9912088597i<br />

c(2) 0 0.1E-12<br />

c(3) 0 0.1E-12<br />

c(5) -1 0.1E-12<br />

c(7) 0 0.1E-12<br />

c(11) −1.11318397521973 0.2E-07<br />

c(13) 0 0.4E-11<br />

c(17) 0 0.1E-10<br />

c(19) −0.17157173156738 0.6E-08<br />

c(23) 0 0.5E-09<br />

c(29) 0.44610100984573 0.1E-07<br />

c(31) −0.16668809652328 0.6E-08<br />

c(37) 0 0.4E-11<br />

c(41) 0.84569376707077 0.1E-07<br />

c(43) 0 0.1E-10<br />

c(47) 0 0.6E-10<br />

c(53) 0 0.1E-09<br />

c(59) −0.59840279817581 0.5E-08<br />

c(61) −1.82808148860931 0.4E-07<br />

c(67) 0 0.1E-08<br />

c(71) 1.29922688007355 0.5E-07<br />

c(73) 0 0.1E-12<br />

c(79) −1.48637688159943 0.2E-07<br />

c(83) 0 0.1E-12<br />

c(89) 1.43317067623138 0.5E-07<br />

c(97) 0 0.1E-12<br />

Some Corresponding Hecke relations (for R = 4.893781291438)<br />

|c(2)c(3) − c(6)| 0.2E-12<br />

|c(93) − c(3)c(31)| 0.3E-11<br />

|c(895) − c(5)c(179)| 0.1E-11<br />

|c(15625) − c(5) 6 | 0.1E-10<br />

||c(5)| − 1| 0.1E-12<br />

|c(195) − c(3)c(5)c(13)| 0.2E-11<br />

|c(9763) − c(13)c(751)| 0.3E-10<br />

|c(72991) − c(47)c(1553)| 0.2E-09<br />

a Since coefficients <strong>of</strong> CM-forms can be computed explicitly (cf. [46, p. 5]), the first<br />

column contains the coefficients computed according to the explicit formula, and<br />

the second column is the difference between this value and the value obtained with<br />

the methods presented in this paper.<br />

b By Proposition 1.2.4(a) we know that |c(5)| = 1.


Table 1.6: Eigenvalues for Γ0(4)<br />

Sine Cosine<br />

R H1 H2 c<br />

3.70330780123 5E-12 7E-12<br />

b 5.41733480684 2E-13 1E-12<br />

6.62042287384 1E-12 5E-13<br />

b 7.22087197596 2E-12 4E-12<br />

b 8.27366588959 4E-12 1E-12<br />

8.52250301688 8E-13 2E-13<br />

a 9.53369526135 3E-12 4E-12<br />

9.93491995937 3E-13 7E-13<br />

b 10.71270690070 2E-13 3E-10<br />

10.76471068319 3E-13 2E-13<br />

b 11.31767970147 4E-13 5E-13<br />

11.97277669404 1E-12 4E-12<br />

a 12.17300832468 2E-12 2E-12<br />

b 12.82198816197 8E-13 4E-13<br />

b 13.31016428347 9E-13 2E-12<br />

b 14.09720373392 9E-12 2E-12<br />

14.1438850658 2E-11 1E-11<br />

a 14.3585095183 2E-11 4E-11<br />

14.97031306857 2E-12 2E-12<br />

b 15.27402248481 2E-12 3E-13<br />

b 15.4428777438 7E-12 6E-11<br />

15.8093206166 2E-11 1E-11<br />

a 16.1380731715 3E-13 2E-11<br />

16.4782589777 1E-11 1E-12<br />

a 16.6442592019 1E-11 8E-12<br />

b 17.3193334036 8E-12 2E-11<br />

b 17.4931127227 2E-11 3E-11<br />

17.6699198653 2E-11 1E-12<br />

R H1 H2 c<br />

5.879354157759 3E-13 6E-13<br />

8.042477591693 7E-12 4E-13<br />

b 8.92287648699 5E-12 2E-12<br />

9.859896162239 7E-13 9E-13<br />

b 10.92039200294 4E-12 3E-12<br />

b 12.09299487508 1E-13 3E-11<br />

12.87761656411 1E-12 2E-13<br />

13.17207496748 5E-12 7E-12<br />

a 13.77975135189 2E-12 1E-10<br />

14.47780273166 1E-13 7E-13<br />

b 14.68501595060 2E-10 7E-12<br />

b 15.31419658418 9E-12 2E-12<br />

15.7434456453 2E-11 6E-11<br />

16.09359237739 7E-12 6E-12<br />

b 16.4041087751 5E-11 7E-11<br />

c 16.94026094636 1E-11 3E-12<br />

17.19512721203 8E-12 2E-12<br />

a 17.73856338106 8E-12 1E-11<br />

b 17.8780025394 2E-11 2E-11<br />

18.24637789203 1E-12 3E-12<br />

18.6207750965 1E-11 2E-11<br />

b 19.12542240816 8E-12 1E-11<br />

a 19.42348147083 9E-12 6E-12<br />

19.69901979586 1E-12 7E-12<br />

19.7785769776 1E-11 7E-11<br />

a Oldform from Γ0(1). This eigenspace is 3-dimensional.<br />

b Oldform from Γ0(2). This eigenspace is 2-dimensional.<br />

c From Proposition 1.2.4(c) we know that |c(4)| = 0 for newforms, and it is easy<br />

to see that in an oldspace f±(z) = f (z) ± ε f (dz), d|4, ε ∈ {±1} (cf. notation<br />

in Remark 1.3.1, and hence c+(n) = c−(n) for (n,4) = 1. Accordingly we set<br />

H2 = |c(4)| for newforms and H2 = |c+(5) − c−(5)| for oldforms.<br />

53


2 Computation <strong>of</strong> <strong>Maass</strong> Waveforms with<br />

Non-trivial Multiplier Systems<br />

2.1 Introduction<br />

In this chapter we will consider computational aspects <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong><br />

on the modular group PSL(2,Z), and the Hecke-congruence subgroups Γ0(N),<br />

non-zero weight and a non-trivial multiplier system.<br />

First we will introduce the basic theory <strong>of</strong> multiplier systems and then we<br />

will introduce the notion <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong> in this context. We will also<br />

provide some details on the different operators that acts on spaces <strong>of</strong> <strong>Maass</strong><br />

<strong>waveforms</strong> with non-trivial multiplier system.<br />

As we will see in section 2.8, we only need to make a few minor changes in<br />

the algorithm we used to compute <strong>Maass</strong> <strong>waveforms</strong> <strong>of</strong> weight 0 (cf. Chapter<br />

1) to adapt it to the case <strong>of</strong> non-zero weight and non-trivial multiplier system.<br />

The biggest obstacle to overcome is the numerical evaluation <strong>of</strong> the Whittaker<br />

function, which is the non-zero weight analogue <strong>of</strong> the K-Bessel function.<br />

The algorithm used to compute this function is presented in Chapter 4.<br />

The last section contains some examples <strong>of</strong> the results that has been obtained<br />

with the described methods. To my knowledge, the only previous computational<br />

investigation <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong> with non-zero weight was performed<br />

by Mühlenbruch, [70], but his method has less accuracy and a smaller<br />

range <strong>of</strong> applicability than the methods to be presented in this chapter.<br />

2.2 Multiplier Systems<br />

2.2.1 Introduction<br />

We will give a brief introduction to multiplier systems, for more extensive<br />

treatments see [39, pp. 331-338], [77], or [84, pp. 70-87].<br />

Let Γ be a Fuchsian group and m be even. Classically, a function ϕ, meromorphic<br />

on the upper half-plane H, which satisfies<br />

ϕ(Az) = ΘA(z;m)ϕ(z) = (cz + d) m <br />

a b<br />

ϕ(z), ∀A = ∈ Γ, (2.1)<br />

c d<br />

55


is called an automorphic form <strong>of</strong> weight m. The function<br />

ΘA(z;m) = (cz + d) m <br />

, A =<br />

a<br />

c<br />

b<br />

d<br />

∈ Γ<br />

is said to be an automorphy factor on Γ. The classical theory <strong>of</strong> automorphic<br />

forms is well-known; for instance, if m = 2, then the automorphic forms can be<br />

identified with the meromorphic differential forms <strong>of</strong> degree 1 on the orbifold<br />

(classical Riemann surface) Γ\H.<br />

We observe that, for even m, the number (cz + d) m is uniquely defined and<br />

the automorphy factor ΘA(z;m) in (2.1) clearly satisfies<br />

ΘA(Bz;m)ΘB(z;m) = ΘAB(z;m). (*)<br />

To generalize these notions to arbitrary real m, there needs to be a choice<br />

<strong>of</strong> branch <strong>of</strong> the argument, and to make certain everything is well-defined, we<br />

have to introduce the notion <strong>of</strong> a multiplier system.<br />

Definition 2.2.1. We will always use the principal branch <strong>of</strong> the argument,<br />

−π < Arg(w) ≤ π. Define<br />

jA(z;m) = e imArg(cz+d) m <br />

(cz + d)m cz + d 2 a b<br />

= =<br />

, A = ∈ SL(2,R).<br />

|cz + d| m cz + d<br />

c d<br />

To adapt the relation (*), we also write<br />

σm(A,B) = jA(Bz;m) jB(z;m) jAB(z;m) −1 . (2.2)<br />

It is clear that for integer m, σm(A,B) = 1, but it can also be shown (cf. [77,<br />

§2, pp. 42–50]) that the only values which σm can take are {1,e ±2πim }.<br />

It is customary to write σm(A,B) = exp[2πimw(A,B)] as in [84, eq. (3.2.1)].<br />

Let Γ be the inverse image <strong>of</strong> Γ under the natural homomorphism SL(2,R) →<br />

PSL(2,R). (Note that this forces −I ∈ Γ.)<br />

Definition 2.2.2. A multiplier system v <strong>of</strong> weight m on Γ is a function v : Γ →<br />

S 1 (where S 1 ⊂ C is the unit circle), such that<br />

• v(−I) = e −πim , and<br />

• v(AB) = σm(A,B)v(A)v(B), ∀A,B ∈ Γ.<br />

Observe that v can be regarded equally well as a multiplier system <strong>of</strong> any<br />

weight m ′ ≡ m mod 2. The question <strong>of</strong> whether there exist multiplier systems<br />

<strong>of</strong> a given weight and on a given group is most easily answered by the<br />

following proposition (cf. [39, prop. 2.1, p. 333]).<br />

Proposition 2.2.1. Given v : Γ → S 1 and m ∈ R. The following are equivalent:<br />

56


• v(T ) is a multiplier system <strong>of</strong> weight m on Γ.<br />

• There exists a function ϕ ≡ 0 on H which is either C∞ or meromorphic<br />

such that<br />

ϕ(T z) = v(T )ϕ(z)(cz + d) m <br />

, ∀T =<br />

a<br />

c<br />

b<br />

d<br />

∈ Γ.<br />

A function ϕ as above which satisfies<br />

ϕ(T z) = v(T )ϕ(z)(cz + d) m , ∀T =<br />

<br />

a b<br />

c d<br />

<br />

∈ Γ,<br />

is called an automorphic form <strong>of</strong> weight m and multiplier system v on Γ.<br />

Definition 2.2.3. Given a multiplier system v on Γ and an element α ∈ GL(2,R)<br />

we define a multiplier system, vα , on the group α−1Γα by<br />

<br />

αAα−1 ,α<br />

v α (A) = v αAα<br />

−1 σm<br />

σm (α,A)<br />

, A ∈ α −1 Γα.<br />

That this indeed gives a multiplier system on α −1 Γα is shown in [63, p.<br />

138].<br />

Using Prop. 2.2.1, we will construct the two most popular multiplier systems<br />

in the following sections. Compare: [39, pp. 334-337].<br />

2.2.2 The η multiplier system<br />

The η-function<br />

The Dedekind η−function is a holomorphic function on H, defined by<br />

<br />

z<br />

∞<br />

η(z) = e (1 − e(nz)),<br />

24<br />

∏<br />

n=1<br />

where e(x) = e2πix . Note that η can be expressed by the following Fourier<br />

series:<br />

η(z) =<br />

<br />

∑ cne n +<br />

n≥0<br />

1<br />

<br />

z ,<br />

24<br />

⎧<br />

⎪⎨ 1, n = 6m<br />

cn =<br />

⎪⎩<br />

2 ± m,<br />

−1, n = 6m2 0,<br />

± 5m + 1,<br />

otherwise.<br />

It is clear from the definition that η(z) = 0 for z ∈ H and that, for each k ∈ R,<br />

η 2k can be defined as a holomorphic function on H. (Cf. [84, p. 205].) It is<br />

57


easily seen that η 2k has a Fourier series <strong>of</strong> the form<br />

η(z) 2k = e<br />

kz<br />

12<br />

<br />

∑ Pn(k)e(nz),<br />

n≥0<br />

where Pn(k) is defined by the relation<br />

j<br />

∏ 1 − q<br />

j≥1<br />

2k = ∑ Pj(k)q<br />

j≥0<br />

j .<br />

Using the expansion 1 − q j2k 2k n<br />

= ∑n≥0 n (−1) q jn , which is valid since<br />

|q| = |e(z)| < 1 for z ∈ H, the first few values <strong>of</strong> Pj(k) can be evaluated to:<br />

P0(k) = 1, P1(k) = −2k, P2(k) = 2k 2 − 3k, P3(k) = − 4<br />

3 k3 + 6k 2 − 8<br />

3 k,<br />

etc. (see a few more at [70, p. 18]). Note that η(z) 24 is the famous Ramanujan<br />

function ∆(z). Cf. [84, pp. 196-197].<br />

The multiplier system<br />

It can be proved (cf. Thm 3.1 and Thm 3.4 [3, p. 48 and p. 52]) that η satisfies<br />

the following functional equations<br />

<br />

−1<br />

η = (−iz)<br />

z<br />

1 2 η (z), and in general,<br />

η(γz) = vη(γ)(cz + d) 1 <br />

2 η(z), ∀γ =<br />

a<br />

c<br />

b<br />

d<br />

∈ SL(2,Z). (2.3)<br />

This functional equation expresses the fact that η is a SL(2,Z)−automorphic<br />

form <strong>of</strong> weight 1<br />

2 and multiplier system given by vη. Accordingly the function<br />

η 2k is an SL(2,Z)-automorphic form <strong>of</strong> weight k and multiplier system given<br />

by v 2k<br />

η , and we can use η 2k in the context <strong>of</strong> Proposition 2.2.1 to assure the<br />

existence <strong>of</strong> the multiplier system, vη,k = v 2k<br />

η , <strong>of</strong> weight k on SL(2,Z) (and<br />

any <strong>of</strong> its subgroups, e.g. Γ0(N)).<br />

We have the following explicit formula for v = v2k η :<br />

1<br />

2πi logv<br />

<br />

kb<br />

a b<br />

=<br />

12 ,<br />

c d k<br />

a = d = 1,c = 0,<br />

a+d−3c<br />

12c − s(d,c) , c > 0,<br />

(2.4)<br />

and <br />

v<br />

−a<br />

−c<br />

−b<br />

−d<br />

= e kπi <br />

v<br />

a<br />

c<br />

b<br />

d<br />

, for c > 0.<br />

58


Here s(d,c) is the Dedekind sum,<br />

s(d,c) =<br />

c−1<br />

∑<br />

n=1<br />

n<br />

c<br />

<br />

dn<br />

,<br />

c<br />

where ((x)) is the saw-tooth function<br />

<br />

x − ⌊x⌋ −<br />

((x)) =<br />

1<br />

2 , if x /∈ Z,<br />

0, if x ∈ Z,<br />

and ⌊x⌋ is the greatest integer less than or equal to x. Note that if x is not<br />

an integer, then ⌊−x⌋ = −⌊x⌋ − 1 so ((−x)) = −((x)), and hence s(−d,c) =<br />

−s(d,c) if gcd(d,c) = 1. One can also express s(d,c) as (cf. [3, p. 61])<br />

s(d,c) = ∑<br />

n mod c<br />

n<br />

c<br />

dn<br />

c<br />

<br />

.<br />

Remark 2.2.1. It is known that, for each k ∈ R, there exist exactly 6 different<br />

multiplier systems <strong>of</strong> weight k on PSL(2,Z) (cf. [84, §3.4, pp. 83, 206 ]<br />

or [63, thm. 19, p. 132]). We will denote these by v (r)<br />

η,k<br />

= v2(k+r) η , where r ∈<br />

{0,2,4,6,8,10}. Compare [84, eq. (6.4.7)]; one knows, <strong>of</strong> course, that v 24<br />

η = 1.<br />

When dealing with the modular group and weight k, it is sufficient to con-<br />

sider only the multiplier system v (0)<br />

later, cf. Section 2.4.5).<br />

2.2.3 The θ Multiplier System<br />

η,k = vη,k = v 2k<br />

η (for reasons to be discussed<br />

On any subgroup <strong>of</strong> PSL(2,Z), we can always use the η-multiplier system, but<br />

in general, if the subgroup has more generators and relations, there are also<br />

other multiplier systems available. In particular, on Γ0(4), there is a multiplier<br />

system <strong>of</strong> weight 1<br />

2 which is interesting from an arithmetical point <strong>of</strong> view.<br />

It is well-known (cf. [90] or [51]) that the Jacobi theta function<br />

θ(z) =<br />

∞<br />

∑<br />

−∞<br />

e(n 2 z), z ∈ H,<br />

is automorphic on Γ0(4) with weight k = 1<br />

2 and can be used to define a multiplier<br />

system on Γ0(4). Using the Poisson summation formula one can prove<br />

(cf. [34, pp. 72–75] or [51, pp. 167–168]) that the θ-function satisfies:<br />

<br />

−1<br />

θ = (−iz)<br />

2z<br />

1 <br />

z<br />

<br />

2 θ , (2.5)<br />

2<br />

59


and one can also prove that (cf. [51, thm. 10.10, p. 177] or [90, p. 447]):<br />

θ(Az) = vθ (cz + d) 1 <br />

2 θ(z), A =<br />

a<br />

c<br />

b<br />

d<br />

∈ Γ0(4), (2.6)<br />

where the multiplier vθ can be expressed explicitly as<br />

<br />

c<br />

<br />

vθ = ¯εd ,<br />

d<br />

where εd = 1 if d ≡ 1 mod 4 and εd = i if d ≡ −1 mod 4, and <br />

c<br />

d denotes the<br />

extended quadratic residue symbol defined as the traditional Jacobi symbol if<br />

0 < d ≡ 1 mod 2 and extended by<br />

<br />

c<br />

<br />

=<br />

d<br />

c<br />

<br />

c<br />

, c = 0,<br />

|c| −d<br />

and<br />

<br />

0 1 if d = ±1,<br />

=<br />

d 0 otherwise.<br />

For the sake <strong>of</strong> completeness we also use the traditional Kronecker extension,<br />

i.e. we define c<br />

2<br />

<br />

=<br />

2<br />

c<br />

<br />

. (2.7)<br />

Observe that since vθ (−I) = e−πik (which follows from the choice <strong>of</strong> argument,<br />

−π < Arg(z) ≤ π), we can always suppose that c > 0 when we evaluate<br />

the symbol <br />

c<br />

d .<br />

One can verify that the symbol <br />

·<br />

· defined as above satisfies reciprocity<br />

relations similar to the usual ones:<br />

Proposition 2.2.2. Suppose that c,d ∈ Z are odd and c = 0. Then we have:<br />

<br />

−1<br />

d<br />

<br />

c<br />

<br />

d<br />

d−1<br />

= (−1)<br />

( 2 )<br />

,<br />

=<br />

<br />

d<br />

c (−1)<br />

( d−1<br />

2 )( c−1<br />

2 )<br />

, d, or c > 0,<br />

− d<br />

c<br />

Worth noting here is that for odd c,d we have<br />

60<br />

d−1<br />

(−1)<br />

( 2 )( c−1<br />

2 )<br />

=<br />

(−1) ( d+1<br />

2 )( c+1<br />

2 ) , d, and c < 0.<br />

<br />

−1, c ≡ d ≡ 3 mod 4,<br />

1, otherwise,


and<br />

d+1<br />

(−1)<br />

( 2 )( c+1<br />

2 )<br />

=<br />

<br />

−1, c ≡ d ≡ 1 mod 4,<br />

1, otherwise.<br />

Remark 2.2.2. Relations (2.5) and (2.6) can also be proved using the corresponding<br />

relations (2.3) for η and the following relation between the η and<br />

the θ functions (cf. [51, p. 177] or [55, thm. 12, p. 46])<br />

<br />

θ(z) = η2 z+1<br />

2<br />

η(z + 1)<br />

. (2.8)<br />

Remark 2.2.3. By (2.8) it is clear that we can define the function θ(z) 2k<br />

uniquely in H for all k ∈ R, and hence we can also define the multiplier system<br />

v 2k<br />

θ <strong>of</strong> weight k on Γ0(4). In this chapter, however, we will consider only half<br />

integral weights for the theta multiplier system.<br />

2.2.4 Further properties <strong>of</strong> the multiplier systems<br />

We are going to consider the eta multiplier system on subgroups <strong>of</strong> the modular<br />

group. It is possible to express the eta multiplier system without involving<br />

a Dedekind sum, i.e. we have the following formula from Knopp [55, p. 51]:<br />

<br />

a b<br />

<br />

c<br />

<br />

1 2<br />

vη<br />

= e (a + d)c − bd c − 1 + 3d − 3 − 3cd<br />

c d d 24<br />

<br />

,<br />

c<br />

d<br />

for c > 0, even. (Note that Knopp’s symbol <br />

c<br />

d<br />

if c = 0).<br />

∗<br />

is equal to our symbol<br />

Remark 2.2.4. Using this formula and the corresponding formula for vθ ([55,<br />

p. 51, thm. 3]) it is easy to see that the eta and theta multipliers coincide on the<br />

principal congruence subgroup Γ(24). Note that [Knopp’s θ](z) =[our θ] <br />

z<br />

2 ,<br />

and that<br />

<br />

<br />

b<br />

a b<br />

a<br />

[Knopp’s vθ ]<br />

= [our vθ ] 2 , when<br />

c d<br />

2c d<br />

<br />

a b<br />

1 0<br />

≡<br />

mod 2.<br />

c d<br />

0 1<br />

Another property <strong>of</strong> the eta multiplier system is that for integral weights it<br />

reduces to a Dirichlet character on certain subgroups.<br />

61


Proposition 2.2.3. The eta multiplier system at weight 1, v2 η, agrees with the<br />

odd Dirichlet character d−1<br />

−1<br />

d = (−1) 2 modulo 4 on the group<br />

<br />

<br />

Γ12 = A ∈ Γ(1)|A ≡<br />

∗<br />

0<br />

0<br />

∗<br />

mod 12 ⊂ Γ0(12)<br />

(and hence it is trivial on Γ(24)).<br />

Pro<strong>of</strong>. Suppose that<br />

A =<br />

<br />

a 12b<br />

12c d<br />

<br />

∈ Γ12.<br />

Then, for c > 0 we have (note that (12,d) = 1)<br />

v 2 η(A) =<br />

<br />

1 2<br />

e (a + d)12c − 12bd 144c − 1 + 3d − 3 − 36cd<br />

12<br />

<br />

=<br />

<br />

1<br />

d − 1<br />

e [3d − 3] = e = i<br />

12 4<br />

d−1 = (−1) d−1<br />

=<br />

2<br />

<br />

−1<br />

.<br />

d<br />

If c = 0 we use (2.4) and see that if a = d = 1 then<br />

and if a = d = −1 then<br />

Hence, in all cases, if c ≥ 0 we have<br />

2.3 <strong>Maass</strong> Waveforms<br />

v 2 η(A) = 1,<br />

v 2 η(A) = −1.<br />

v 2 η(A) =<br />

<br />

−1<br />

.<br />

d<br />

The slash-operator f |A(z) = f (Az) can be extended to an operator <strong>of</strong> weight k<br />

as:<br />

f |[A,k](z) = f (Az) jA(z;k) −1 ,<br />

62


and the natural analog <strong>of</strong> the Laplace-Beltrami operator, ∆, which is invariant<br />

under this action is the weight-k Laplacian:<br />

∆k = ∆ − iyk ∂<br />

<br />

∂ 2 ∂ 2<br />

= y2 +<br />

∂x ∂x2 ∂y2 <br />

− iyk ∂<br />

∂x .<br />

If Γ is a Fuchsian group we define the space M(Γ,v,k,λ) consisting <strong>of</strong> <strong>Maass</strong><br />

<strong>waveforms</strong> on Γ, <strong>of</strong> weight k, multiplier system v and eigenvalue λ, as the<br />

space <strong>of</strong> functions which satisfy the following conditions:<br />

1) f |[A,k](z) = v(A) f (z), ∀A =<br />

2) ∆k f + λ f = 0, and<br />

3)<br />

<br />

F | f |2 dµ < ∞.<br />

Observe that condition 1) is equivalent to<br />

1’) f (Az) = v(A) jA(z;k) f (z), ∀A =<br />

<br />

a b<br />

c d<br />

<br />

<br />

a b<br />

c d<br />

∈ Γ,<br />

<br />

∈ Γ.<br />

For purposes <strong>of</strong> the computational work to be described in this chapter, we<br />

shall be content to restrict ourselves to cases where λ > 1<br />

4 . (Cf. also here para.<br />

4 <strong>of</strong> sect. 2.9.1 below.)<br />

Instead <strong>of</strong> the Bessel equation in the case <strong>of</strong> weight 0, condition 2) above<br />

gives us the Whittaker equation, and using the method <strong>of</strong> separation <strong>of</strong> variables<br />

gives us Whittaker functions instead <strong>of</strong> the K-Bessel functions at weight<br />

0 (for complete details see [39, chap. 9]). Since f (x+iy) is no longer periodic<br />

in x, but instead satisfies f (z + 1) = v(S) f (z) = e(α) f (z), with α ∈ [0,1), the<br />

Fourier series <strong>of</strong> f can ([39, pp. 26, 348, 420(19)]) be written as<br />

f (z) =<br />

∞<br />

∑<br />

−∞<br />

n+α=0<br />

c(n)<br />

W k sgn(n+α)<br />

|n + α| 2 ,iR (4π|n + α|y)e((n + α)x), (2.9)<br />

where Wl,µ(x) is the Whittaker function in standard notation (cf. [29, vol. I, p.<br />

264]) and R is the usual spectral parameter, λ = 1<br />

4 + R2 . One notes here that<br />

Wl,µ(x) = e− x 2 x (µ+ 1 <br />

2) 1<br />

x<br />

Ψ µ + .<br />

2 − l;2µ + 1;x and that W0,µ(x) = π− 1 2 x 1 2 Kµ 2<br />

For k = 0, the expansion above thus reduces to usual Fourier expansion with<br />

2y 1 2 KiR(2π|n + α|y) as in [39, p. 26, prop. 4.12].<br />

As in the case <strong>of</strong> trivial multiplier, we define functions f j related to f at<br />

each cusp, p j, <strong>of</strong> Γ by f j(z) = f |[σ j,k](z) = jσ j (z;k)−1 f (σ jz). Here σ j is the<br />

63


cusp normalizing map <strong>of</strong> the cusp p j (cf. Chapter 1, Definition 1.1.4), and<br />

using σ jSσ −1<br />

j = S j together with (2.2) it easy to see that<br />

f j(z + 1) = v(S j) f j(z) = e(α j) f j(z),<br />

with α j ∈ [0,1) (cf. [51, p. 41]). Thus the Fourier series <strong>of</strong> f at the cusp j can<br />

be written as<br />

c j(n)<br />

f j(z) = <br />

|n + α j| Wsgn(n+α j) k 2 ,iR(4π|n + α j|y)e((n + α j)x). (2.10)<br />

∞<br />

∑<br />

n=−∞<br />

As in the case <strong>of</strong> weight 0 and Dirichlet character, we say that the cusp number<br />

j is open or singular if α j = 0 and closed if α j = 0. If all cusps <strong>of</strong> Γ are singular<br />

for the multiplier system v we say that v is a singular multiplier system for Γ.<br />

Remark 2.3.1. Observe that, for the eta-multiplier and weight k, we have<br />

α = α1 = k<br />

12 .<br />

2.3.1 Decomposition <strong>of</strong> the discrete spectrum<br />

It is known (see for example [18] or [39]) that closed cusps (i.e. v(S j) = 1)<br />

do not contribute to the continuous spectrum, and if all cusps are closed there<br />

is only the discrete part <strong>of</strong> the spectrum left and this is spanned by the <strong>Maass</strong><br />

<strong>waveforms</strong>. We also know (see [39, p. 385]) that on the modular group with<br />

weight k the least eigenvalue is<br />

λmin = |k|<br />

2<br />

<br />

1 − |k|<br />

<br />

,<br />

2<br />

or larger. In the case <strong>of</strong> PSL(2,Z) and k ≥ 0, F(z) = y k 2 η(z) 2k has eigenvalue<br />

equal to λmin.<br />

In this chapter, any eigenvalues λ ∈ λmin, 1<br />

<br />

4 will be regarded as exceptional.<br />

The non-exceptional eigenvalues thus satisfy 1<br />

4 < λ0 ≤ λ1 ≤ ··· ≤<br />

λn → ∞. One can obtain lower bounds for the eigenvalue λ0 (see [18, p. 183]),<br />

but in light <strong>of</strong> the numerical experiments in Section 2.9.1 they are not very<br />

effective (cf. Figure 2.2).<br />

64


2.4 Operators<br />

2.4.1 Conjugation and reflection<br />

Let J and K denote the conjugation, Kz = z, and reflection, Jz = −z. Then J<br />

and K act as involutions on the space <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong> via the operations<br />

K f = f |K(z) = f (z).<br />

J f = f |J(z) = f (−z).<br />

The action <strong>of</strong> GL(2,R) on H is defined as follows; for g = <br />

a b<br />

c d and z ∈ H<br />

we define<br />

⎧<br />

az+b<br />

⎪⎨ cz+d , if ad − bc > 0,<br />

g(z) =<br />

⎪⎩ az+b<br />

cz+d , if ad − bc < 0,<br />

then as usual we can use the following matrix in GL(2,R) to represent the<br />

operator J :<br />

For T = <br />

a b<br />

c d we define<br />

J =<br />

<br />

T ∗ = JT J −1 =<br />

1 0<br />

0 −1<br />

<br />

<br />

.<br />

a −b<br />

−c d<br />

and then T ∗∗ = T, and T (z) |K = −T ∗ (z), meaning that also −T z = T ∗ (−¯z).<br />

Remark 2.4.1. It is easy to verify that if f ∈ M(Γ0(N),v,k,λ), then the conjugate<br />

K f ∈ M(Γ0(N),v,−k,λ) and J f ∈ M(Γ0(N),v ∗ ,−k,λ), where v∗ is the<br />

multiplier system determined by<br />

v ∗ (T ) = v(T ∗ <br />

) ·<br />

1,<br />

e<br />

c = 0,<br />

πik(1−sgn(d)) <br />

,<br />

, forT =<br />

c = 0,<br />

a<br />

c<br />

b<br />

d<br />

.<br />

Here we used the fact that for c = 0<br />

jT ∗(Jz;k) = eikArg(−c(−¯z)+d) = e ikArg(c¯z+d) = e −ikArg(cz+d) = jT (z;k) −1 ,<br />

while for c = 0 we have<br />

jT ∗ (Jz;k) = eikArg(d) =<br />

<br />

<br />

1 = jT (z;k) −1 , d > 0,<br />

e ikπ = e 2πik jT (z;k) −1 , d < 0.<br />

,<br />

65


What we want to study here is the involution obtained by combining J and<br />

K, i.e.<br />

KJ f (z) = f |JK(z) = f (−¯z).<br />

It is easily seen that if f has Fourier coefficients c j(n), then (by [29, p. 265(8)])<br />

f |JK has Fourier coefficients c j(n), and we thus would like to have f and f |JK<br />

belonging to the same space (i.e. transform according to the same multiplier<br />

system), since then we can assume that the Fourier coefficients are real.<br />

It is clear that if f ∈ M(Γ0(N),v,k,λ), then KJ f ∈ M(Γ0(N),v ∗ ,k,λ) so we<br />

are left to see whether v ∗ = v or not. We consider the two multiplier systems<br />

separately.<br />

Eta multiplier<br />

Suppose that v = vη, T = <br />

a b<br />

c d , and without loss <strong>of</strong> generality we assume<br />

c ≥ 0. We consider the case c > 0 first:<br />

1<br />

2πi lnv(T ∗ ) = 1<br />

2πi lnv<br />

<br />

a<br />

−c<br />

−b<br />

d<br />

= k 1<br />

+<br />

2 2πi lnv<br />

<br />

=<br />

−a b<br />

c −d<br />

k<br />

=<br />

=<br />

<br />

−a − d − 3c<br />

+ k<br />

− s(−d,c)<br />

2 12c<br />

<br />

−a − d + 3c<br />

k<br />

− s(−d,c)<br />

12c<br />

<br />

a + d − 3c<br />

−k<br />

− s(d,c)<br />

12c<br />

= − 1<br />

lnv(T ),<br />

2πi<br />

so v(T ∗ ) = v(T ), and hence v ∗ (T ) = v J (T ) = v(T ). Now, if c = 0, and d > 0<br />

we have<br />

66<br />

v(T ∗ <br />

) =<br />

=<br />

1 −b<br />

v<br />

0 1<br />

<br />

−kb<br />

e<br />

24<br />

= v(T ),


and if c = 0, d < 0, then<br />

v(T ∗ <br />

) =<br />

=<br />

=<br />

−1 −b<br />

v<br />

0 −1<br />

<br />

k 1<br />

e v<br />

2 0<br />

<br />

k kb<br />

e + ,<br />

2 24<br />

<br />

b<br />

1<br />

and<br />

hence v(T ∗ ) = e(−k)v(T ), and thus<br />

Theta multiplier<br />

<br />

v(T ) =<br />

=<br />

=<br />

−1 b<br />

v<br />

0 −1<br />

<br />

k 1<br />

e v<br />

2 0<br />

<br />

k kb<br />

e − ,<br />

2 24<br />

<br />

−b<br />

1<br />

v ∗ (T ) = v(T ∗ )e(k) = v(T ).<br />

In the case <strong>of</strong> the θ-multiplier we have, for c > 0 or c = 0 and d > 0 :<br />

<br />

a −b<br />

vθ,k (T ∗ ) = vθ<br />

= εd 2k<br />

= εd 2k<br />

=<br />

−c d<br />

2k −c<br />

d<br />

2k −1<br />

<br />

c<br />

2k d d<br />

2k , ifd ≡ 1 mod 4,<br />

i2k <br />

c 2k<br />

d , if d ≡ 3 mod 4,<br />

c<br />

d<br />

= vθ (T ),<br />

67


and in the case c = 0, d < 0 we have<br />

<br />

−1 b<br />

and<br />

vθ (T ∗ ) = vθ<br />

0 −1<br />

= ε−1 2k<br />

2k 0<br />

−1<br />

= (−i) 2k<br />

= vθ (T ),<br />

<br />

v ∗ (T ) = vθ (T ∗ )e 2kπi = i 2k (−1) 2k = (−i) 2k = vθ (T ),<br />

hence in all cases, v ∗ θ (T ) = vθ (T ), and we conclude this section with the following<br />

proposition.<br />

Proposition 2.4.1. If v is either the η- or the θ-multiplier system (in the latter<br />

case 4|N) then a basis { f1,..., fm} <strong>of</strong> M(Γ0(N),v,k,λ) can be chosen so that<br />

each f j can be expanded in a Fourier series at ∞ with real coefficients.<br />

Pro<strong>of</strong>. We have seen that for both the theta and the eta multiplier systems the<br />

product KJ is a conjugate-linear involution <strong>of</strong> the space M(Γ0(N),v,k,λ), and<br />

hence we can assume that any f ∈ M(Γ0(N),v,k,λ) is an eigenfunction <strong>of</strong><br />

KJ with eigenvalue ε, where |ε| = 1. Note that if f (z) has a Fourier series<br />

expansion as above, then<br />

f |J|K(z) = f (−¯z)<br />

= ∑<br />

n+α=0<br />

c(n)<br />

W k<br />

|n + α| 2 sgn(n+α),ν (4π|n + α|)e(x(n + α))<br />

hence c(n) = εc(n). Finally we observe that if ε = e iθ we can look at the function<br />

g = e i θ 2 f which then satisfies KJg = e −i θ 2 KJ f = e −i θ 2 e iθ f = g, hence we<br />

can (after proper normalization) assume that the eigenvalue <strong>of</strong> KJ is ε = 1, and<br />

that the Fourier coefficients are real.<br />

Remark 2.4.2. For the sake <strong>of</strong> completeness it should be remarked that in general<br />

one can not simultaneously take Fourier coefficients at cusps other than ∞<br />

to be real (cf. Subsection 2.4.3 and the map ωN, which is a cusp normalizing<br />

map for the cusp at 0 and which has eigenvalues ±i −2k ).<br />

2.4.2 The involution ωN<br />

As usual we define ωNz = −1<br />

Nz , or equivalently ωN =<br />

−1 <br />

0 √N<br />

√ . We know that<br />

N 0<br />

ωN is an involution <strong>of</strong> Γ0(N), i.e. Γ0(N) = ωNΓ0(N)ω −1<br />

n , but the question is<br />

68


how it relates to the weight and multiplier system.<br />

If f ∈ M(Γ0(N),v,k,λ) it is easy to see that f |[k,ωN] ∈ M(Γ0(N),v ωN ,k,λ),<br />

and that vωN <br />

(T ) = v ωNT ω −1<br />

<br />

N . We also have<br />

f |[k,ωN]|[k,ωN](z) = jωN (z;k)−1 f (ωNz) |[k,ωN]<br />

= jωN (z;k)−1 jωN (ωNz;k) −1 f ω 2 Nz <br />

= e −ikArg(√ Nz) e −ikArg(−1/ √ Nz) f (z)<br />

= e −iπk f (z),<br />

and hence if we define the operator τN by<br />

we have that<br />

τN f (z) = e ik π 2 f |[k,ωN](z),<br />

τ 2 N = Id.<br />

To prove that τN is a linear involution, what is left to verify is that v ωN = v,<br />

which we can do in the following two cases:<br />

a) If N = 4, and v = vθ , using θ <br />

−1<br />

1<br />

2z = (−iz) 2 θ <br />

z<br />

2 one can prove<br />

that vω4(A) = v(A).<br />

b) If N = 1, and v = vη, by evaluating v ω1 on the generators <strong>of</strong> Γ0(1),<br />

S = ω1, and T, it is immediate that v ω1<br />

η = vη.<br />

And we conclude:<br />

Proposition 2.4.2. If N = 1 and v = vη, or N = 4 and v = vθ the operator<br />

defined by<br />

τN : M(Γ0(N),v,k,λ) → M(Γ0(N),v,k,λ),<br />

τN f (z) = e ik π 2 f |[k,ωN] = e −ik(Argz− π 2 ) f<br />

is a linear involution, i.e. τ2 N f = f .<br />

<br />

−1<br />

,<br />

Nz<br />

This means that the operator τN is a linear involution <strong>of</strong> M(Γ0(N),v,k,λ),<br />

and hence has eigenvalues ±1. Note that in the case N = 4, if τN f (z) = ± f (z),<br />

then<br />

f2 = f |ωN = e−ik π 2 τN f = ±e −ik π 2 f ,<br />

which in turn means that the Fourier coefficients at the cusp at 0 satisfy<br />

c2(n) = ±e −i π 2 k c1(n) = ±(i) −k c1(n). (2.11)<br />

69


2.4.3 The operator σ4<br />

For N = 4 and θ-multiplier system we define the operator σ4, by<br />

<br />

z + j<br />

.<br />

σ4 f (z) = ∑<br />

j mod 4<br />

f<br />

This can be viewed as a “partial” Hecke operator on Γ0(4). (cf. [53, p.<br />

195, (0.8)] and [56, p. 250]). We have to verify that σ4 is an operator on<br />

M(Γ0(4),vθ ,k,λ). Suppose that A = <br />

a b<br />

4c d ∈ Γ0(4), and with the notation<br />

Vj = <br />

1 j<br />

, we get<br />

0 4<br />

(σ4 f ) |[A,k](z) = jA(z;k) −1 ∑<br />

j mod N<br />

f (VjAz)<br />

Here we used<br />

jB(Vlz;k) = exp<br />

= jA(z;k) −1 ∑<br />

l mod N<br />

f<br />

<br />

4<br />

∗ ∗<br />

16c d − 4lc<br />

<br />

Vlz<br />

{l is uniquely defined mod 4 by al ≡ b + jd mod 4}<br />

= jA(z;k) −1 ∑ ε<br />

l mod N<br />

2k<br />

2k 16c<br />

d<br />

vθ (B) jB(Vlz;k) f (Vlz)<br />

d − 4lc<br />

= ε2k <br />

16c<br />

d ∑<br />

f (Vlz).<br />

d − 4lc<br />

= jA(z;k).<br />

l mod N<br />

<br />

ikArg 16c<br />

z + l<br />

4<br />

<br />

<br />

+ d − 4lc = exp[ikArg(4cz + d)]<br />

Hence, what is left is to check that <br />

16c c c<br />

d−4lc = d−4lc = d .<br />

The sign <strong>of</strong> Nc is the same as the sign <strong>of</strong> N2c and we might assume without<br />

loss <strong>of</strong> generality that c > 0. By virtue <strong>of</strong> (2.7) we may also assume that c is<br />

odd since we can always pull out a factor <strong>of</strong> 2 at the beginning. For the reciprocity<br />

relations applicable to our quadratic residue symbol, see Proposition<br />

2.2.2:<br />

70<br />

<br />

c<br />

d − 4lc<br />

<br />

=<br />

=<br />

=<br />

<br />

d − 4lc<br />

c−1<br />

(−1) ( 2 )(<br />

c<br />

d−4lc−1<br />

2 )<br />

<br />

d<br />

c−1<br />

(−1) ( 2 )(<br />

c<br />

d−1<br />

<br />

2 )−l(c−1)c d<br />

c−1<br />

= (−1) ( 2 )(<br />

c<br />

d−1<br />

2 )<br />

<br />

c<br />

<br />

,<br />

d


(note that c(c − 1) is even). We have shown that<br />

(σ4 f ) |[A,k](z) = ε 2k<br />

d<br />

<br />

c<br />

2k d<br />

∑<br />

l mod N<br />

and thus σ4 : M(Γ0(4),vθ ,k,λ) → M(Γ0(4),vθ ,k,λ).<br />

f (Vlz) = vθ,k(A)σ4 f (z),<br />

Remark 2.4.3. Following Kohnen [56] or Katok-Sarnak [53] we define the<br />

operator L acting on M(Γ0(4),vθ , 1<br />

2 ,λ) by<br />

L = τ4σ4.<br />

It can be shown that this operator is self-adjoint, commutes with ∆ 1 and have<br />

2<br />

the eigenvalues 1 and 1<br />

2 .<br />

2.4.4 <strong>Maass</strong> operators<br />

So far, the operators we have seen act on spaces <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong> <strong>of</strong> a<br />

given weight and multiplier system.<br />

We will show that we may limit the range <strong>of</strong> weights k to investigate to k ∈<br />

[0,6]. For this we will use the <strong>Maass</strong> lowering and raising operators, E ± k ,which<br />

raise or lower the weight <strong>of</strong> a <strong>Maass</strong> waveform by units <strong>of</strong> 2. They are defined<br />

by<br />

E ± k<br />

∂ ∂ k<br />

= iy ± y +<br />

∂x ∂y 2 ,<br />

and using the relation between the Whittaker function and the confluent hypergeometric<br />

function together with the transformation formulas [29, p. 258,<br />

(10)] (see also: [63, p. 183 (middle)] and [64, p. 302 lines -3 and -1]) we see<br />

that for k > 0 (use the abbreviation 4π|n + α|y = Y )<br />

<br />

<br />

(Y )e((n + α)x) = −Wk+2<br />

E + k<br />

E − k<br />

E + k<br />

E − k<br />

W k 2 ,iR<br />

<br />

W k − 2 ,iR (Y )e((n + α)x)<br />

<br />

=<br />

2<br />

<br />

k(k + 2)<br />

4<br />

,iR (Y )e((n + α)x) if n + α > 0,<br />

+ 1<br />

<br />

+ R2<br />

4<br />

W k+2 − 2 ,iR (Y )e((n + α)x)<br />

<br />

<br />

if n + α < 0,<br />

W k − 2 ,iR (Y )e((n + α)x) =W k−2 − 2 ,iR (Y )e((n + α)x)<br />

<br />

<br />

k(k − 2)<br />

W k<br />

2 ,iR (Y )e((n + α)x) = − +<br />

4<br />

if n + α < 0, and<br />

1<br />

<br />

+ R2<br />

4<br />

ifn + α > 0.<br />

Wk−2<br />

2 ,iR (Y )e((n + α)x)<br />

To verify that they respect the weight k-action we rewrite them in the follow-<br />

71


ing form<br />

E + k<br />

E − k<br />

and also use the formulation:<br />

∂ k<br />

= (z − z) +<br />

∂z 2 ,<br />

∂ k<br />

= −(z − z) +<br />

∂z 2 ,<br />

e −iArg(cz+d) =<br />

1<br />

cz + d 2<br />

.<br />

cz + d<br />

It is then easy to see (cf. [63, p. 178]) that<br />

E + <br />

k<br />

f|[k,A](z) <br />

= (z − z) ∂<br />

<br />

k<br />

+<br />

∂z 2<br />

k <br />

cz + d 2<br />

f (Az)<br />

cz + d<br />

k+2<br />

−k cz + d 2 c<br />

=(z − z)<br />

f (Az)<br />

2 cz + d (cz + d)<br />

k<br />

cz + d 2 1<br />

+<br />

cz + d (cz + d) 2 f ′ <br />

(Az) + k<br />

k<br />

cz + d 2<br />

f (Az)<br />

2 cz + d<br />

k <br />

cz + d 2 k c(z − z)<br />

(z − z)<br />

=<br />

1 − f (Az) +<br />

cz + d 2 (cz + d)<br />

k+2<br />

cz + d 2<br />

=<br />

cz + d<br />

(z − z)<br />

|cz + d| 2 f ′ (Az) + k<br />

<br />

f (Az)<br />

2<br />

k+2<br />

cz + d 2<br />

=<br />

cz + d<br />

<br />

′ k<br />

Az − Az f (Az) + f (Az)<br />

2<br />

= E + <br />

k f |[k+2,A] .<br />

(cz + d) 2 f ′ (Az)<br />

It is known that E ± k maps M(Γ0(N),v,k,λ) into M(Γ0(N),v,k±2,λ), and that<br />

the composition<br />

E ∓ k±2 E± k : M(Γ0(N),v,k,λ) ↦→ M(Γ0(N),v,k,λ)<br />

is just multiplication by a constant, which is readily seen to be nonzero anytime<br />

λ > 1<br />

4 . Hence E± k acts bijectively on the spaces corresponding to nonexceptional<br />

eigenvalues, i.e. they are always bijections for λ > 1<br />

4 .<br />

2.4.5 <strong>Maass</strong> operators and the symmetry about k = 6<br />

First <strong>of</strong> all, observe that E ± k<br />

only change the weight and not the multiplier<br />

system v, but in view <strong>of</strong> the remark after Def. 2.2.2 it is clear that v = v 2k<br />

η is<br />

72


also a multiplier system <strong>of</strong> weight k + r for any r ∈ 2Z, and with the notation<br />

it is clear that<br />

v (r)<br />

η,k<br />

and<br />

= v2(k+r)<br />

η<br />

E + k<br />

E − k<br />

(r)<br />

(r−2)<br />

: M(Γ0(N),v η,k ,k,λ) → M(Γ0(N),v η,k+2 ,k + 2,λ),<br />

(r)<br />

(r+2)<br />

: M(Γ0(N),v η,k ,k,λ) → M(Γ0(N),v η,k−2 ,k − 2,λ).<br />

Our main purpose is to investigate the locations <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong> on<br />

the modular group when the weight and multiplier system are varied. That<br />

is, we would like to investigate the space M(Γ0(1),v (r)<br />

η,k ,k,λ) for all k ∈ R<br />

and r ∈ {0,2,4,6,8,10}. To reduce the problems we use the fact that v (r)<br />

η,k are<br />

all 24th-roots <strong>of</strong> unity. Suppose that λ > 1<br />

4 so the the lowering and raising<br />

operators act bijectively, then using the notation (k,v (r)<br />

η,k ) to denote that space<br />

M(Γ0(1),v (r)<br />

η,k ,k,λ) we have:<br />

• A composition <strong>of</strong> <strong>Maass</strong> operators which raises or lowers the weight by 12<br />

will preserve the multiplier system. Hence it is no restriction to assume<br />

that k ∈ [0,12].<br />

• Conjugation maps f ∈ (k,v (r)<br />

η,k ) to K f ∈ (−k,v(−r) η,−k ) = (−k,v(12−r)<br />

η,−k ), which<br />

can in turn be mapped to a function g ∈ (12 − k,v (12−r)<br />

η,12−k ) by raising operators,<br />

so there is no restriction to assume k ∈ [0,6] and r ∈ {0,2,4,6,8,10}.<br />

• But by using the raising operator repeatedly we can map v (0)<br />

η,k to any v(r)<br />

η,k<br />

and hence it is no restriction to also assume r = 0.<br />

We are thus justified in our choice <strong>of</strong> restricting the investigation to the spaces<br />

M(Γ0(1),v,k,λ) for k ∈ [0,6] and v = v (0)<br />

η,k<br />

2.4.6 Hecke operators<br />

= v2k<br />

η .<br />

We know that Hecke operators play an important role in the understanding <strong>of</strong><br />

the theory <strong>of</strong> both modular forms and <strong>Maass</strong> <strong>waveforms</strong> at integer weights.<br />

For general real weights, the Hecke operators are not important, but we will<br />

begin with the general definition anyway, and then we will consider the case<br />

<strong>of</strong> weight 1<br />

2<br />

and the theta multiplier system, where the action <strong>of</strong> the Hecke<br />

operators T p 2 are well understood.<br />

This discussion is based on [93], and see also [103].<br />

Let Γ ⊂ PSL(2,R) be c<strong>of</strong>inite and suppose v : Γ→ S 1 is a multiplier system<br />

<strong>of</strong> weight k ∈ R. The commensurator <strong>of</strong> Γ, comm(Γ), in PSL(2,R) is defined<br />

as<br />

comm(Γ) = α ∈ PSL(2,R)|αΓα −1 ∩ Γ has finite index in Γ and αΓα −1 .<br />

73


We know that the Hecke operators are associated with the members <strong>of</strong> the<br />

commensurator, or actually with the double cosets, ΓαΓ, for α ∈ comm(Γ).<br />

Fix α ∈ comm(Γ). It can be shown that if the multiplier system v satisfies<br />

v(g) = v α (g), ∀g ∈ Γ ∩ α −1 Γα, (2.12)<br />

then we can define an associated multiplier system, vα, on the double coset:<br />

by setting<br />

vα : ΓαΓ → S 1 ,<br />

vα (g1αg2) = σk (g1α,g2)σk (g1,α)v(g1)v(g2),<br />

for all g1,g2 ∈ Γ. It might be the case that there exists an associated multiplier<br />

system <strong>of</strong> the multiplier system W = χv, where χ is a character on ΓαΓ, even<br />

though there does not exist an associated multiplier system <strong>of</strong> v itself. Suppose<br />

that vα exists as above, and that ΓαΓ = ∪ d i=1 Γαi. We then define the Hecke<br />

operator T v α,k : M Γ,v,k,λ → M Γ,v,k,λ by<br />

T v α,k f (z) =<br />

=<br />

d<br />

∑ vα (αi) f |[αi,k](z)<br />

i=1<br />

d<br />

∑ vα (αi) jαi<br />

i=1<br />

(z;k)−1 f (αiz).<br />

When Γ is a congruence subgroup, one usually constructs Hecke operators<br />

corresponding to positive integers n, in which case α = <br />

1 0<br />

0 n and αi = <br />

a b<br />

0 d<br />

with ad = n, d > 0, b mod d. The Hecke operator T v<br />

n,k is then defined by<br />

T v<br />

n,k f (z) = n k 2 −1 <br />

d−1<br />

<br />

a b az + b<br />

∑ ∑ vα<br />

f .<br />

0 d d<br />

ad=n b=0<br />

The Theta multiplier<br />

Consider the case Γ = Γ0(4), k = 1<br />

2 and v = vθ . Let n be a positive integer and<br />

α = <br />

1 0<br />

0 n . Then g ∈ Γ0(4) ∩ α−1 <br />

a nb<br />

Γ0 (4)α can be written g = c/n d , and<br />

αgα−1 = <br />

a b , with ad − bc = 1 and c ≡ 0 mod 4n. It is easy to verify that<br />

c d<br />

vα (g) = v(αgα−1 ), and<br />

c<br />

n<br />

v(g) = εd<br />

d<br />

<br />

= εd<br />

c<br />

d<br />

<br />

n<br />

−1 = v(αgα<br />

d<br />

−1 )<br />

<br />

n<br />

−1 ,<br />

d<br />

and hence vα (g) = v(g) if and only if <br />

n<br />

d = 1. Since (2.12) implies this condition<br />

for all odd integers d it is clear that the extension vα exists if and only<br />

if n is a perfect square.<br />

74


Suppose now for simplicity that n = p2 , with p a prime number = 2. It<br />

is easy to verify that in this case the associated multiplier system is given<br />

by vα (g1αg2) = vα (g1)vα (g2), and the different coset representatives are<br />

1 b<br />

given by 0 p2 <br />

,b = 0,..., p2 <br />

− 1, p2 0<br />

p b<br />

, and<br />

0 1<br />

0 p , b = 0,..., p − 1. By<br />

factorization <strong>of</strong> these representatives we see that<br />

<br />

p<br />

vα<br />

2 <br />

b<br />

1 0<br />

= vα<br />

0 1<br />

0 p2 <br />

1 b<br />

= 1,<br />

0 1<br />

<br />

p2 <br />

0<br />

p2 <br />

−t 1 0 p2 <br />

d t<br />

vα<br />

vα<br />

<br />

and hence<br />

0 1<br />

= vα<br />

4 d 0 p2 −4 1<br />

= p 2 d + 4t = 1 <br />

p<br />

= v<br />

2 <br />

−t<br />

4 d<br />

= εd = 1,<br />

<br />

p b<br />

0 p<br />

= vα<br />

T vθ<br />

p2 , 1 f (z)<br />

2<br />

= p − 3 =<br />

d−1<br />

2 ∑ ∑ vα<br />

ad=p2 b=0<br />

p − 3 <br />

p2−1 2 ∑ vα<br />

b=0<br />

vα<br />

= p − 3 2<br />

= p − 3 2<br />

1 0<br />

4xp 1<br />

<br />

1 0<br />

0 p 2<br />

<br />

p b<br />

−4x D<br />

= {pD + 4bx = 1}<br />

<br />

p b<br />

−4x<br />

= v<br />

= εD<br />

−4x D<br />

D<br />

<br />

−4x −4x −4x<br />

= εp<br />

= εp<br />

pD p<br />

p<br />

<br />

4xb −b −b b<br />

= εp<br />

= εp = εp ,<br />

p p p p<br />

<br />

<br />

a b az + b<br />

f<br />

0 d d<br />

1 b<br />

0 p2 <br />

z + b<br />

f<br />

p2 <br />

+<br />

<br />

p2 <br />

0<br />

f<br />

0 1<br />

p 2 z <br />

p−1<br />

<br />

p b<br />

+ ∑ vα<br />

f z +<br />

b=0 0 p<br />

b<br />

<br />

p<br />

<br />

<br />

p2−1 <br />

z + b<br />

∑ f<br />

b=0 p2 <br />

+ f p 2 z p−1 <br />

b<br />

+ ∑ εp f z +<br />

b=0 p<br />

b<br />

<br />

p<br />

<br />

<br />

z + b<br />

f<br />

p2 <br />

+ f p 2 z p−1<br />

<br />

b<br />

+ εp ∑ f z +<br />

p<br />

b<br />

<br />

p<br />

<br />

.<br />

p 2 −1<br />

∑<br />

b=0<br />

b=0<br />

75


This agrees with the corresponding results in [90, pp. 450–451]. To see that<br />

the operator obtained is the same as in [53] (up to a constant multiple) we look<br />

at the action on the Fourier coefficients a(n) <strong>of</strong> f ∈ M Γ0(4),vθ , 1<br />

2 ,λ , i.e.<br />

with the abbreviation W(x) = W1<br />

4 sgn(x),iR (4π|x|), and<br />

we have<br />

and note that<br />

and<br />

T vθ<br />

p2 , 1 f (z) = p<br />

2<br />

− 3 2<br />

p−1<br />

∑<br />

b=0<br />

<br />

b<br />

e<br />

p<br />

f (z) = ∑ n=0<br />

+ ∑ n=0<br />

+ εp<br />

= p − 3 2<br />

+ ∑ n=0<br />

+ ∑ n=0<br />

p 2 −1<br />

∑<br />

b=0<br />

a(n)<br />

W(ny)e(nx),<br />

|n|<br />

∑<br />

n=0<br />

<br />

a(n) ny<br />

W<br />

|n| p2 <br />

n(x + b)<br />

e<br />

p2 <br />

a(n)<br />

W(np<br />

|n| 2 y)e(np 2 x)<br />

p−1<br />

∑<br />

b=0<br />

p2−1 ∑ e<br />

b=0<br />

<br />

nb<br />

p<br />

(cf. [61, IV, §3]) which gives us 1<br />

<br />

<br />

b a(n)<br />

p ∑ W(ny)e nx +<br />

n=0 |n| nb<br />

<br />

p<br />

<br />

<br />

a(n)<br />

∑ W<br />

n=0 |n|<br />

<br />

ny<br />

p2 <br />

e<br />

a(n)<br />

W(np<br />

|n| 2 y)e(np 2 x)<br />

a(n)<br />

W(ny)e(nx)εp<br />

|n|<br />

<br />

nx<br />

p2 p2−1 p−1<br />

∑<br />

b=0<br />

<br />

nb<br />

p2 <br />

p<br />

=<br />

2 , p2 |n,<br />

0, p2 ∤ n,<br />

∑ e<br />

b=0<br />

<br />

b<br />

e<br />

p<br />

<br />

nb<br />

p2 <br />

<br />

nb<br />

p<br />

<br />

,<br />

p−1<br />

<br />

n nb nb<br />

=<br />

p ∑ e<br />

b=0 p p<br />

= √ <br />

n<br />

p<br />

p<br />

<br />

1, p ≡ 1 mod 4,<br />

i, p ≡ −1 mod 4,<br />

= √ <br />

n<br />

p εp,<br />

p<br />

1 For p = 2, we use a definition analogous to [90, p. 450] and derive an expansion analogous to<br />

[90, thm. 1.7].<br />

76


T vθ<br />

p2 , 1 f (z) =<br />

2<br />

1<br />

√<br />

p ∑<br />

n=0<br />

1<br />

<br />

|n|<br />

<br />

a(p 2 <br />

n<br />

n) + a<br />

p2 <br />

+ p − 1 2<br />

<br />

n<br />

a(n) W(ny)e(nx).<br />

p<br />

It is easy to see that our Hecke operator T vθ<br />

p2 , 1 =<br />

2<br />

1 √<br />

pTp2, where Tp2 is the corresponding<br />

Hecke operator defined in [53, p. 199]. Observe that our Fourier<br />

coefficients a(n) = √ n×Katok-Sarnak’s Fourier coefficients b(n).<br />

As usual, we can look at Hecke eigenforms in M Γ0(4),vθ , 1<br />

2 ,λ , and one<br />

can prove (cf. [90, p. 453]) that if f = ∑ a(n)<br />

√ W(ny)e(nx) is an eigenfunction<br />

n<br />

<strong>of</strong> T vθ<br />

p2 , 1 for all primes p and if t is square free, then<br />

2<br />

a(tm 2 )a(tn 2 ) = a(t)a(tm 2 n 2 ), for(m,n) = 1.<br />

2.4.7 Lifts at weight 1 and Fourier coefficients<br />

As a general rule, in cases with nontrivial multiplier system, it is hard to see<br />

precisely which Hecke operators (if any) are present and what the corresponding<br />

multiplicative relations are. In this section we discuss a simple reduction<br />

trick which allows one to get some multiplicative relations in a case where<br />

k = 1 and α = 0 (cf. (2.9)).<br />

To this end, consider any <strong>Maass</strong> <strong>waveforms</strong> f ∈ M(Γ0(1),v2 η,1,λ). This<br />

implies that, for ad − 144bc = 1, we have (see Prop. 2.2.3)<br />

f (Az) = jA(z;1)v 2 <br />

−1<br />

η (A) f (z) = jA(z;1) f (z),<br />

d<br />

where A = <br />

a 12b<br />

12c d , i.e. f ∈ M(Γ12, <br />

−1<br />

d ,1,λ). We now follow a standard<br />

idea to “lift” f to a function <strong>of</strong> weight 1 on Γ0(144) with a Dirichlet character.<br />

Indeed, set<br />

g(z) = f (12z).<br />

Then, with<br />

A =<br />

<br />

a b<br />

144c d<br />

<br />

∈ Γ0(144),<br />

77


we see that<br />

<br />

az + b 12az + 12b<br />

g(Az) = g<br />

= f<br />

144cz + d 144cz + d<br />

<br />

a(12z) + 12b<br />

= f<br />

12c(12z) + d<br />

= e iArg(144cz+d)<br />

<br />

−1<br />

f (12z)<br />

= e iArg(144cz+d)<br />

d<br />

−1<br />

d<br />

<br />

g(z),<br />

i.e. g ∈ M(Γ0(144), <br />

−1<br />

d ,1,λ). The advantage <strong>of</strong> this transformation is that<br />

the theory <strong>of</strong> Hecke operators is well understood in the case <strong>of</strong> a Hecke congruence<br />

subgroup, Γ0(N), together with a Dirichlet character χ. In fact, the<br />

Hecke operators on M(Γ0(144), χ,1,λ) are the same as at weight 0:<br />

<br />

az + b<br />

Tn f = ∑ ∑ χ(d) f ,<br />

d<br />

ad=n<br />

d>0<br />

b mod d<br />

and the multiplicativity relation is the usual (cf. [69, p. 37 (para 3) and §4.5.]):<br />

TnTm = ∑<br />

d|(m,n)<br />

χ(d)dTmn .<br />

d2 We will now consider how the existence <strong>of</strong> these Hecke operators gives rise<br />

to relations between the Fourier coefficients <strong>of</strong> f . Assume that we have<br />

∞ a(n)<br />

f (z) = ∑ n + 1<br />

<br />

<br />

<br />

4π <br />

1 <br />

n + <br />

12<br />

y<br />

<br />

e n + 1<br />

<br />

x ,<br />

12<br />

then<br />

f (12z) =<br />

n=−∞<br />

∞<br />

∑<br />

n=−∞<br />

= √ 12<br />

a(n)<br />

n <br />

1 +<br />

∞<br />

∑<br />

n=−∞<br />

<br />

12<br />

W1 2 sgn(n),iR<br />

<br />

12<br />

W1 2 sgn(n),iR<br />

and hence the Fourier expansion <strong>of</strong> 1<br />

√12 g is<br />

78<br />

∑<br />

n≡1 mod 12<br />

<br />

<br />

4π <br />

1 <br />

n + <br />

12<br />

12y<br />

<br />

e n + 1<br />

<br />

12x<br />

12<br />

a(n)<br />

W1<br />

|12n + 1| 2 sgn(n),iR (4π |12n + 1|y)e((12n + 1)x),<br />

c(n)<br />

W1<br />

|n| 2 sgn(n),iR (4π|n|y)e(nx),


where<br />

<br />

n − 1<br />

c(n) = a .<br />

12<br />

The multiplicativity relations for the Hecke operators imply similar relations<br />

for the normalized coefficients c(n)/c(1) :<br />

<br />

−1<br />

<br />

mn<br />

c<br />

d d2 <br />

, (2.13)<br />

c(m)c(n) = c(1) ∑<br />

0


and that if g ∈ M(Γ0(144), <br />

−1<br />

d ,1,λ) satisfy Λg = εg with ε = ±1, and have<br />

Fourier coefficients c(n) as above, then<br />

c(−n) = ±iRc(n).<br />

Using this in combination with (2.13) we get the following relation<br />

c(−n)c(−m) = −R 2 c(1) ∑<br />

00<br />

−1<br />

d<br />

<br />

c<br />

mn<br />

d 2<br />

valid for −n ≡ −m ≡ 1 mod 12 (which implies mn ≡ 1 mod 12). If we change<br />

to a’s we get that for any positive integers m,n we have<br />

<br />

−1 (12m − 1)(12n − 1) − d2 a<br />

d<br />

12d2 <br />

.<br />

And to get a simpler relation (that is easy to verify numerically) we set<br />

m = −1, and we get two cases. If (12n − 1,11) = 1 then<br />

and if (12n − 1,11) = 11 then<br />

<br />

,<br />

a(−n) = −R2a(0) a(11n − 1), (2.17)<br />

a(−1)<br />

a(−n) = −R2 a(0)<br />

a(−1)<br />

<br />

a(11n − 1) − a<br />

n − 1<br />

11<br />

<br />

. (2.18)<br />

We have not seen relations <strong>of</strong> type (2.16)-(2.18) earlier in the literature. For<br />

some numerical examples, see Tables 2.2 and 2.3.<br />

Remark 2.4.4. It is clear that not all forms on M(Γ0(144), <br />

−1<br />

d ,1,λ) correspond<br />

to forms on M(Γ0(1),v2 η,1,λ) (look at the Weyl’s law for the two cases).<br />

From experiments we see that M(Γ0(1),v2 η,1,λ) seems to contain cosine forms<br />

<strong>of</strong> CM-type as well as other “generic” forms. See Table 2.4.<br />

2.5 Oldforms<br />

Clearly the usual f ∈ M(Γ0(M),v,k,λ) ⇒ f ∈ M(Γ0(N),v,k,λ) holds whenever<br />

M|N. The obvious notation is to call such a function f an oldform on<br />

M(Γ0(N),v,k,λ). We will look at one example <strong>of</strong> an oldform.<br />

80


2.5.1 Γ0(N) with N prime and η-multiplier<br />

The cusps are p1 = ∞ and p2 = 0, with cusp normalizing maps and stabilizers<br />

(i.e. generators <strong>of</strong> the stabilizer subgroups) <strong>of</strong> σ1 = Id, S1 = S and<br />

<br />

σ2 =<br />

0<br />

N<br />

−1<br />

0<br />

S2 = σ2Sσ −1<br />

2 =<br />

<br />

1<br />

−N<br />

0<br />

1<br />

.<br />

The shifts are α1 = k<br />

and let T = 0 −1<br />

1 0<br />

12 = Nα1.<br />

<br />

Suppose that f ∈ M(Γ0(1),v,k,λ),<br />

and dN = <br />

N 0<br />

0 1 so σ2 = T dN then<br />

12 and α2 = Nk<br />

f2(z) = f |σ2 (z) = jσ2 (z;k)−1 f (σ2z)<br />

= jσ2 (z;k)−1 f (T dNz)<br />

= jσ2 (z;k)−1 jT (dNz;k)v(T ) f (dNz)<br />

= e −ikArgz e ikArgz v(T ) f (dNz)<br />

= e −ikπ v(T ) f (Nz)<br />

= e −i π 2 k f (Nz) = (−i) k f (Nz).<br />

In terms <strong>of</strong> the Fourier series this means that<br />

f2(z) =<br />

c2(n)<br />

<br />

|n + α2| Wsgn(n+α2) k 2 ,µ(4π|n + α2|y)e((n + α2)x)<br />

f (Nz) =<br />

=<br />

=<br />

=<br />

∞<br />

∑<br />

−∞<br />

∞<br />

∑<br />

−∞<br />

∞<br />

∑<br />

−∞<br />

∞<br />

∑<br />

−∞<br />

∞<br />

∑<br />

m=−∞<br />

c1(n)<br />

<br />

|n + α1| Wsgn(n+α1) k 2 ,µ(4π|n + α1|Ny)e((n + α1)Nx)<br />

c1(n)<br />

<br />

|n + α2<br />

N |<br />

W α<br />

sgn(n+ 2<br />

N ) k 2 ,µ(4π|Nn + α2|y)e((Nn + α2)x)<br />

√<br />

Nc1(n)<br />

<br />

|Nn + α2| Wsgn(Nn+α2) k 2 ,µ(4π|Nn + α2|y)e((Nn + α2)x)<br />

√ <br />

m Nc1 N <br />

|m + α2| Wsgn(m+α2) k 2 ,µ(4π|m + α2|y)e((m + α2)x),<br />

where as usual c(r) = 0 if r ∈ Z, and since f2(z) = (−i) k f (Nz) we must have<br />

the following relation:<br />

c2(Nn) = (−i) k√ Nc1(n).<br />

81


2.6 The Eisenstein series<br />

In case there is a cusp p j at which the multiplier system is singular (i.e.<br />

v(S j) = 1) we have a continuous spectrum: [ 1<br />

4 ,∞) (with multiplicity equal<br />

to the number <strong>of</strong> singular cusps), and in general we can not say much about<br />

the embedded discrete spectrum in [ 1<br />

4 ,∞).<br />

Singular cusps can appear either because <strong>of</strong> the multiplier system, e.g. the<br />

θ-multiplier system, or because the “twist” α j is an integer (remember v(S j) =<br />

e 2πiα j ), e.g. on the modular group with the η-multiplier system and weight<br />

k = 12l.<br />

Remember that <strong>Maass</strong> <strong>waveforms</strong> are part <strong>of</strong> the discrete spectrum, but as<br />

we continuously “turn <strong>of</strong>f” the multiplier system, i.e. for v = v 2k<br />

η we let k → 0,<br />

the continuous spectrum will emerge in the limit. For this reason we want to<br />

review some details concerning the Eisenstein series on the modular group,<br />

i.e. for the rest <strong>of</strong> this section Γ = PSL(2,Z).<br />

2.6.1 Weight 0<br />

At weight 0 and singular character χ, the continuous spectrum <strong>of</strong> ∆ = ∆0 is the<br />

interval [ 1<br />

4 ,∞) and the eigenpacket is given by the Eisenstein series E(z;s; χ)<br />

defined by<br />

E(z;s; χ) = ∑<br />

T ∈Γ∞\Γ<br />

χ(T −1 )(ℑ(T z)) s ,<br />

where Γ∞ = [S]. (See [39, p. 65]). From proposition 6 <strong>of</strong> the same page we get<br />

E(z;s; χ) = y s + ϕ(s)y 1−s ∑ ϕm(s)<br />

m=0<br />

√ yK 1 s− (2π|m|y)e(mx)<br />

2<br />

here if W∞ = <br />

a b<br />

c d are the double coset representatives, i.e. Γ = ∑Γ∞W∞Γ∞<br />

then<br />

ϕ(s) = √ π Γs − 1<br />

<br />

2<br />

∞ )<br />

Γ(s)<br />

1<br />

,<br />

|c| 2s<br />

and<br />

82<br />

ϕm(s) = 2πs |m| s− 1 2<br />

Γ(s)<br />

∑ χ(W<br />

W∞∈Γ∞<br />

−1<br />

∑ χ(W<br />

W∞∈Γ∞<br />

−1<br />

∞ ) 1<br />

|c|<br />

2s e<br />

<br />

md<br />

.<br />

c


For trivial character χ = 1 we get (see page 76 <strong>of</strong> the above)<br />

ϕ(s) = √ π Γs − 1<br />

<br />

2 ζ (2s − 1)<br />

Γ(s) ζ (2s)<br />

ϕm(s) = 2πs |m| s− 1 2<br />

Γ(s)<br />

σ1−2s(|m|)<br />

.<br />

ζ (2s)<br />

Hence we can see that for s = 1<br />

2 + iR , the n−th Fourier coefficient <strong>of</strong> E(z;s)<br />

is given by<br />

<br />

1<br />

c(n) = ϕ + iR = K · |n|<br />

2 iR σ−2iR(|n|), (2.19)<br />

= K · |n| iR ∑<br />

d||n|,d>0<br />

d −2iR<br />

where K = K(R) is a constant dependent on R. First we observe the simple<br />

formula for prime coefficients (p is positive prime)<br />

c(p) = K · p iR (1 + p −2iR )<br />

= 2K · cos(Rln p)<br />

so we can compute quotients <strong>of</strong> c(p) and compare this with corresponding<br />

quotients for our forms <strong>of</strong> weight = 0.<br />

2.7 Half integer weight<br />

We know that the θ−function is an automorphic form (not a cusp form) <strong>of</strong><br />

weight 1<br />

2 on Γ0(4), hence we consider Γ0(4) together with the θ−multiplier<br />

system (cf. section 2.2.3).<br />

2.7.1 The Shimura correspondence<br />

Introduction – The Holomorphic Case<br />

We will consider the Shimura correspondence only in the particular case <strong>of</strong><br />

trivial character and level a square-free multiple <strong>of</strong> 4. Let Sk(M) denote the<br />

space <strong>of</strong> holomorphic cusp forms <strong>of</strong> weight k ∈ Z (and trivial multiplier) on<br />

Γ0(M), and let S 1 k+ (4N) denote the space <strong>of</strong> holomorphic cusp forms <strong>of</strong><br />

2<br />

weight k + 1<br />

2 , k ∈ Z and multiplier vθ , on Γ0(4N). The Shimura correspondence<br />

is a correspondence between the space S 1 k+ (4N) and spaces S2k(N<br />

2<br />

′ )<br />

for certain integers N ′ |4N (e.g. N ′ = 2N or N).<br />

The map from S 1 k+ to S2k was first constructed by Shimura [90] and later<br />

2<br />

was constructed by Shintani [92]. The for-<br />

an adjoint map from S2k to S k+ 1 2<br />

83


mer uses a Dirichlet-series and the latter uses an integral against a thetafunction.<br />

Both these maps commute with the Hecke operators that are acting<br />

on S2k(N) and S 1 k+ (4N) respectively. Kohnen [56, 57] proved that for<br />

2<br />

N odd and square-free, the correspondence is a bijection between the newforms<br />

on S2k(N) and a certain subspace, V + ⊆ S 1 k+ (4N). The subspace V<br />

2<br />

+<br />

is composed <strong>of</strong> Hecke eigenfunctions whose Fourier coefficients, c(n), satisfy<br />

certain vanishing properties; namely, c(n) = 0 for n ≡ 2,3 mod 4. (Here we<br />

shall mainly be interested in N = 1.)<br />

Extension to the Non-holomorphic case<br />

The extension <strong>of</strong> the Shimura correspondence and Kohnen’s result to spaces<br />

<strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong> has been investigated by e.g. Sarnak [87], Hejhal [37],<br />

Duke [28], Katok-Sarnak [53], Kojima [58, 59], Biró [13] and Arakawa [4].<br />

It is not clear exactly what the range and domain <strong>of</strong> the correspondence are<br />

in the <strong>Maass</strong> case.<br />

Define the two spaces U + ⊆ M(Γ0(1),0,2R) and V + ⊆ M Γ0(4),vθ , 1<br />

2 ,R<br />

<strong>of</strong> Hecke normalized <strong>Maass</strong> <strong>waveforms</strong>. Here U + consists <strong>of</strong> the even (cosine)<br />

Hecke eigenforms on Γ0(1) and V + consists <strong>of</strong> the Hecke eigenforms (<strong>of</strong> the<br />

) which are invariant under the operator L = τ4σ4 (cf.<br />

Hecke operators T vθ<br />

p 2 , 1 2<br />

Remark 2.4.3). As in the holomorphic case one can show that the condition<br />

L f = f is equivalent to the statement that the Fourier coefficients <strong>of</strong> f satisfy<br />

c(n) = 0 for n ≡ 2,3 mod 4 (cf. [56, prop. 2].<br />

From [53, prop. 4.1 and 2.3] we know that to each φ ∈V + there corresponds<br />

a lift f = F(φ), which is either ≡ 0 or in U + . (The “0” occurs if and only if<br />

the first Fourier coefficient <strong>of</strong> φ is zero.) Conversely it is also true that for<br />

each f ∈ U + there exists φ ∈ V + such that f = F(φ).<br />

By combining what is known from the holomorphic and non-holomorphic<br />

cases, we would also expect to find a correspondence between even newforms<br />

on Γ0(2) <strong>of</strong> weight 0 and spectral parameter 2R and certain (other, non V + )<br />

half integer forms on Γ0(4) with spectral parameter R, and indeed we found<br />

such a correspondence in our numerical investigations which we present in<br />

Subsection 2.9.4.<br />

Fourier coefficients<br />

We will now see how to relate the Fourier coefficients <strong>of</strong> the half integral and<br />

corresponding integral weight form.<br />

There are formulas in [53] and [58] connecting the Fourier coefficients <strong>of</strong><br />

half integral and integral weight, but due to differences in normalizations it<br />

might be clearer to go back to Shimura ([90]) and start there.<br />

Suppose first that we have two Hecke eigenforms: f and φ, where f (z) ∈<br />

M(Γ0(2),0,2R) + (the space <strong>of</strong> even newforms) and φ(z) ∈ M Γ0(4),vθ , 1<br />

2 ,R .<br />

84


Further suppose that their respective Fourier series expansions with respect to<br />

the cusp at ∞ are<br />

and<br />

f (z) =<br />

φ(z) =<br />

∞<br />

∑<br />

n=−∞<br />

n=0<br />

∞<br />

∑<br />

n=−∞<br />

n=0<br />

A(n) √ yK2iR (2π|n|y)e(nx),<br />

a(n)<br />

W1<br />

|n| 4 sgn(n),iR (4π|n|y)e(nx),<br />

where we use the normalization a(1) = 1 and A(1) = 1.<br />

By [53, eq. (4.2)] and the analogy with what happens in the holomorphic<br />

case (cf. [90, pp. 453, 458]), when f is the lift <strong>of</strong> φ, one has the following:<br />

Natural Conjecture. Let f ∈ M(Γ0(2),0,2R) + be the Shimura lift <strong>of</strong> φ ∈<br />

M Γ0(4),vθ , 1<br />

2 ,R . Suppose that f has Fourier coefficients {A(n)} and that φ<br />

has Fourier coefficients {a(n)}(with respect to ∞). Let t ∈ Z + be square free<br />

such that a(t) = 0. Then, for each n ∈ Z + , we have:<br />

A(n) = ∑<br />

kd=n<br />

k>0<br />

χ ′ t (k)<br />

√ k<br />

where χ ′ t is the quadratic residue symbol t<br />

·<br />

448(prop. 1.3), 458 (main theorem), 474 (line -11)].)<br />

Outline <strong>of</strong> pro<strong>of</strong>. Identify coefficients in the relation<br />

∞<br />

∑ A(n)n<br />

n=1<br />

−s = c · L<br />

<br />

s + 1<br />

2 , χ′ t<br />

a td2 , (2.20)<br />

a(t)<br />

considered mod 4t. (Cf. [90, pp.<br />

∞<br />

∑ a<br />

n=1<br />

tn 2 n −s , (2.21)<br />

which can be seen as the <strong>Maass</strong> waveform equivalent <strong>of</strong> [90, p. 458] and [72,<br />

p. 159], or alternatively as an extension (to newforms) <strong>of</strong> [53, eqs. (0.15) and<br />

(4.2)].<br />

To obtain this relation, it seems that it is possible to follow the construction<br />

in [53], but using the theta function from Example 3 in [72, p. 151–153] instead<br />

<strong>of</strong> θ(z,g). In our setting this amounts to twisting the theta function by the<br />

character <br />

t<br />

· , and apart from different transformation rules all calculations in<br />

[53] should go through.<br />

The main difference in the argument is that instead <strong>of</strong> starting with a function<br />

φ(z) on Γ0(4), one starts with φ(tz) on Γ0(4t). Cf. [90, p. 474]. Also see<br />

[54, prop. 4.1 and thm. 5.1(2)]<br />

85


In the earlier (oldform) setup <strong>of</strong> [53], (2.21) should hold with χ ′ t considered<br />

mod t. Likewise for (2.20).<br />

This point requires one clarification, however. Namely: since weight 0<br />

oldforms on Γ0(2) occur in pairs, the Shintani map can be expected to give<br />

rise to a corresponding pair <strong>of</strong> weight 1<br />

2 forms on Γ0(4). In a suitable basis,<br />

one <strong>of</strong> these should lie in the space V + ; the other should not but should still<br />

for p = 2. Cf. [69, lemmas 4.6.2<br />

have the same Hecke eigenvalues under T vθ<br />

p 2 , 1 2<br />

and 4.6.10]. In light <strong>of</strong> this, the second form can still be expected to satisfy<br />

(2.20), but only for (n,4) = 1.<br />

The preceding (partly conjectural) statements are backed by strong numerical<br />

evidence. Part <strong>of</strong> this evidence can be found in Table 2.8. Further conjectures<br />

(backed by numerics) concerning the Shimura correspondence can be<br />

found in Section 2.9.4.<br />

See [19, §4.1] and [21, p. 633], and [91, pp. 502 (bottom) – 504 (top)] for<br />

some additional perspectives.<br />

2.8 Some <strong>Computational</strong> Remarks<br />

A priori there are basically only two modifications that need to be done in<br />

order to make the weight zero algorithm work in the general case:<br />

• the K-Bessel function needs to be replaced with the Whittaker function,<br />

and<br />

• the automorphy condition needs to incorporate the multiplier system.<br />

The first modification, although trivial in theory is the most complex in terms<br />

<strong>of</strong> the numerics involved. There was no efficient algorithm for the Whittaker<br />

function in the literature and we had to develop a new algorithm. We used the<br />

integral representation (cf. [39, p. 375 (top)])<br />

<br />

2x ∞<br />

Wk,iR(2x) = e<br />

π<br />

−xcosht <br />

Ψ<br />

0<br />

−k; 1<br />

;x(1 + cosht)<br />

2<br />

<br />

cosh[iRt]dt,<br />

where Ψ is a confluent hypergeometric function together with a stationary<br />

phase method to develop a robust and efficient algorithm. This algorithm is<br />

in essence similar to Hejhal’s algorithm (cf. [42]) for the K-Bessel function,<br />

KiR(x), and the generalization <strong>of</strong> it made by Avelin (cf. [9]) to a complex<br />

argument, Ks(x), s ∈ C. The details <strong>of</strong> this algorithm are described in Chapter<br />

4.<br />

To modify the automorphy condition is trivial. Observe that f j(z) = f |σ (z) = j<br />

jσ j (z;k)−1 f (σ jz), and with z∗ −1<br />

j = σI( j) UI( j)Tjσ jz (cf. p. 35) the automorphy<br />

condition (1.19) now becomes:<br />

86


f j(z) = jσ j (z;k)−1 f (σ jz) = jσ j (z;k)−1 f (T −1 −1<br />

j UI( j) σI( j)z ∗ j)<br />

= jσ j (z;k)−1v(T −1 −1<br />

j UI( j) ,k) jT −1<br />

j U−1 (σI( j)z<br />

I( j)<br />

∗ j;k) f (σI( j)z ∗ j)<br />

= jσ j (z;k)−1 j −1<br />

Tj U−1 (σI( j)z<br />

I( j)<br />

∗ j;k) jσ (z I( j) ∗ j;k)v(T −1 −1<br />

j UI( j) ,k) fI( j)(z ∗ j).<br />

The entire setup <strong>of</strong> the Phase 1 algorithm (i.e. locating eigenvalues and<br />

computing the first few Fourier coefficients) goes through exactly as in the<br />

case <strong>of</strong> weight 0 (cf. Section 1.3.2) with some trivial modifications. For the<br />

sake <strong>of</strong> completeness we will give an outline <strong>of</strong> the modified algorithm.<br />

Consider a <strong>Maass</strong> waveform, f , <strong>of</strong> weight k and multiplier system v. At<br />

each cusp it has a Fourier series<br />

fi(z) =<br />

∞<br />

∑<br />

−∞<br />

n+α i =0<br />

ci(n)<br />

|n + αi| W sgn(n+αi) k 2 ,iR(4π|n + αi|y)e((n + αi)x),<br />

and since the Whittaker function is ultimately exponentially decaying, given<br />

an ε > 0, there exists a constant (depending on y and ε), M(y), such that<br />

fi(z) = ˆfi(z) + [[ε]],<br />

where we use [[ε]] to denote a constant <strong>of</strong> magnitude less than ε. The truncated<br />

Fourier series<br />

ˆfi(z) =<br />

M(Y )<br />

∑<br />

−M(Y )<br />

n+α i =0<br />

ci(n)<br />

<br />

|n + αi| Wsgn(n+αi) k 2 ,iR(4π|n + αi|y)e((n + αi)x),<br />

is now viewed as a discrete Fourier transform, and if we take the inverse<br />

transform over the horocyclic points: zm = xm + iY , 1 − Q ≤ m ≤ Q, where<br />

− m) we get:<br />

xm = 1 1<br />

2Q ( 2<br />

1<br />

2Q<br />

= 1<br />

2Q<br />

=<br />

Q<br />

∑<br />

1−Q<br />

Q<br />

∑<br />

1−Q<br />

M(Y )<br />

∑<br />

−M(Y )<br />

l+α i =0<br />

ˆfi(zm)e(−(n + α j)xm) =<br />

M(Y )<br />

∑<br />

−M(Y )<br />

l+α i =0<br />

ci(l)<br />

<br />

|l + αi| Wsgn(l+αi) k 2 ,iR(4π|l + αi|Y )e((l + αi)xm − (n + α j)xn)<br />

ci(l)<br />

<br />

|l + αi| Wsgn(l+αi) k 2 ,iR (4π|l + αi|Y ) 1<br />

2Q<br />

= ci(n)<br />

|n + α j| W sgn(n+α j) k 2 ,iR (4π|n + α j|Y ),<br />

Q<br />

∑<br />

1−Q<br />

e((l + αi)xm − (n + α j)xn)<br />

87


where we used the fact that<br />

Q<br />

1<br />

2Q ∑ e((l + αi)xm − (n + α j)xn)<br />

1−Q<br />

=<br />

<br />

(l + αi − n − α j)<br />

e<br />

×<br />

4Q<br />

1<br />

2Q ×<br />

Q <br />

× ∑ e −(l + αi − n − α j) m<br />

<br />

2Q<br />

[since α ∈ [0,1)] = δnl.<br />

1−Q<br />

= δ{l + αi = n + α j}<br />

Now we also have fi(zm) = χmi fI(m,i)(z∗ mi ), where<br />

χmi = jσi (zm;k) −1 j −1<br />

Ti U−1<br />

I(m,i) σI(m,i) (z∗mi;k)w(T −1<br />

i<br />

Hence<br />

where<br />

i j<br />

Vnl =<br />

= 1<br />

2Q<br />

= 1<br />

2Q<br />

=<br />

U −1<br />

I(m,i) ,σI(m,i))v(T −1<br />

i<br />

ci(n)<br />

<br />

|n + αi| Wsgn(n+αi) k 2 ,iR(4π|n + αi|y)<br />

κ<br />

∑<br />

j=1<br />

We then define ˜V<br />

i j<br />

nl by<br />

˜V<br />

i j i j<br />

nl = V<br />

Q<br />

∑<br />

1−Q<br />

fi(zm)e(−(n + αi)xm) + [[ε]]<br />

Q<br />

∑ χmi fI(m,i)(z 1−Q<br />

∗ mi)e(−(n + αi)xm) + [[ε]]<br />

M0<br />

∑<br />

l=−M0<br />

n+α i =0<br />

1 1<br />

<br />

|l + α j| 2Q<br />

c j(l)V<br />

i j<br />

nl<br />

+ 2[[ε]],<br />

U −1<br />

Q<br />

∑ χmiWsgn(l+α j)<br />

1−Q<br />

I(m,i)= j<br />

k 2 ,iR(4π|l + α j|y ∗ m j) ×<br />

×e (l + α j)x ∗ <br />

m j e(−(n + αi)xm).<br />

nl − δnlδ ji<br />

1<br />

<br />

|n + αi| Wsgn(n+αi) k 2 ,iR(4π|n + αi|Y ).<br />

If we now neglect the error ε we have a linear system<br />

I(m,i) ,k).<br />

CV = 0, (*)<br />

which must be satisfied by the Fourier coefficients <strong>of</strong> f . Here V is the matrix<br />

˜V<br />

i j<br />

nl and C is the vector ci(n), both depending on R and Y . The basic idea<br />

88


<strong>of</strong> Phase 1 is now to locate eigenvalues R together with sets <strong>of</strong> corresponding<br />

Fourier coefficients, ci(n), by solving (*) repeatedly for different R’s, and<br />

seeking those values <strong>of</strong> R for which the solution vectors C = C(R,Y ) are stable<br />

under changes <strong>of</strong> Y . That is, if C(R,Y1) ≈ C(R,Y2) (for Y1and Y2 < Y0,<br />

for some suitable constant Y0) we take it as an indication that the R is close<br />

to an eigenvalue and that the components <strong>of</strong> C are close to the corresponding<br />

Fourier coefficients. For more details and justifications see Section 1.3.2.<br />

We note that the Phase 2 algorithm <strong>of</strong> Section 1.3.3 (i.e. computation <strong>of</strong><br />

a larger set <strong>of</strong> Fourier coefficients) also generalizes to non zero weight in a<br />

similar manner.<br />

2.9 Numerical results<br />

The experimental excursions have been directed towards three essentially different<br />

subjects, but, in each, we have worked in an exploratory spirit.<br />

• First we tried to get an over-all picture <strong>of</strong> the distribution <strong>of</strong> small to middlerange<br />

sized eigenvalues on PSL(2,Z) (and the eta multiplier) for “large”<br />

weights, e.g. R ∈ [0,14] and k ∈ [0.1,6].<br />

• Second, we continuously “turned <strong>of</strong>f” the multiplier system v2k η on PSL(2,Z)<br />

by letting the weight k → 0 and studied the varying distribution <strong>of</strong> eigenvalues,<br />

Nk(T ), as well as the formation <strong>of</strong> the continuous spectrum.<br />

• There are some cases where arithmeticity plays a role even in nonzero<br />

weight. We studied the Shimura correspondence between weight zero<br />

forms on Γ0(2) and weight one half forms on Γ0(4). And we also studied<br />

weight one forms on PSL(2,Z) which correspond to Hecke eigenforms<br />

on Γ0(144) with a Dirichlet character.<br />

Further data than that which is presented in this section is available from the<br />

author upon request.<br />

2.9.1 Varying weight<br />

The first experiment considered was to tabulate the first few eigenvalues (up<br />

to R = 14) for PSL(2,Z) and the multiplier system v2k η , <strong>of</strong> weight k ∈ (0,6]<br />

(cf. Section 2.4.5). We made the computations for k ∈ [0.1, 6] with a grid size<br />

<strong>of</strong> 0.01, and the results are presented in Figure 2.1. We stress here, that the<br />

arcs in Figure 2.1 terminate at k = 0.1; it is not excluded that they might go<br />

lower.<br />

For some examples <strong>of</strong> eigenvalues for “large” weights see Table 2.1. This<br />

data should be compared with data obtained by Mühlenbruch, [70, p. 160],<br />

who used a completely different method (with much less accuracy). We note<br />

here that as R increases, the negative Fourier coefficients seem to grow rapidly<br />

89


Figure 2.1: Section <strong>of</strong> eigenvalues with 0 < R ≤ 14, and weight 0.1 ≤ k ≤ 6.<br />

R<br />

14<br />

12<br />

10<br />

8<br />

6<br />

4<br />

2<br />

0<br />

0.1 1 2 3 4 5 6<br />

Weight k<br />

in magnitude (as compared to the positive ones, with the normalization c(1) =<br />

1) for “large” weights.<br />

We believe that we have found all eigenvalues with (R,k) ∈ [0,14]×[0.1, 6].<br />

This belief is founded on the “continuity” <strong>of</strong> the resulting graphs R j(k) (cf. Figure<br />

2.1), where R j(k) is the j-th eigenvalue at weight k, considered as a function<br />

<strong>of</strong> k. By standard results (e.g. [17, p. 149]) R j(k) should be a real analytic<br />

function in this interval.<br />

Remember that, for k ≥ 0, the smallest eigenvalue, λmin, corresponds to<br />

the function F(z) = y k 2 η(z) 2k . A lower bound for the second smallest eigenvalue,<br />

λ0(k), is discussed in [17, p. 183]. Bruggeman finds two such bounds,<br />

both positive, which he calls µ0(k) and µ1(k) (µ1 is better than µ0 in a certain<br />

interval I ⊂ [0,2].) Figure 2.2 shows a comparison between the R-values corresponding<br />

to these bounds and the smallest experimentally found eigenvalues<br />

in the interval k ∈ [0.1, 6]; we see that Bruggeman’s bounds can be improved<br />

quite a bit.<br />

2.9.2 Small weights<br />

The investigation <strong>of</strong> eigenvalues for small weights k has been done in the<br />

interval R ∈ [0,20], and k = 10 − j , where 1 ≤ j ≤ 12, i.e. k ∈ [10 −12 ,0.1]. We<br />

believe that most cusp forms were found.<br />

Let us use the notation<br />

λ j(k)<br />

for the j-th discrete eigenvalue at weight k, and φ j(k) for the corresponding<br />

cusp form. It is then a basic fact that λ j(k) depends continuously on k, but<br />

it can also be shown that for k ∈ (0,12) λ j is even real analytic in k. That is,<br />

φ j(k) belong to an “analytic family” in the terminology <strong>of</strong> Bruggeman (see<br />

90


Figure 2.2: Comparison with the theoretical lower bounds in k ∈ [0.1,6].<br />

R<br />

4.0<br />

3.5<br />

3.0<br />

2.5<br />

2.0<br />

1.5<br />

1.0<br />

0.5<br />

µ1(k)<br />

µ0(k)<br />

0.0<br />

0.1 1 2 3 4 5 6<br />

Weight k<br />

[17, 18]). In connection with this, it should also be noted that our experiments<br />

support the statement in Observation 71 in [70] that the first few cusp forms<br />

at weight 0 do not belong to an analytic family <strong>of</strong> cusp forms defined in the<br />

interval (−12,12). Indeed we find that in the range considered we actually<br />

seem to have λ j(k) → 0 as k → 0 which is consistent with Bruggeman’s result<br />

(cf. case ii) b) in Prop. 2.17 in [17, p. 149]).<br />

Our experiments indicate that for fixed small k, the “generic” cusp forms<br />

φ j(k) can be divided into two classes:<br />

• C(k), and<br />

• E(k).<br />

Here C(k) consists <strong>of</strong> functions φ j(k) such that λ j(k) is close to an eigenvalue<br />

<strong>of</strong> a cusp form at weight 0 and the Fourier coefficients are close to the<br />

corresponding coefficients <strong>of</strong> the weight-zero cusp form.<br />

E(k) on the other hand consists <strong>of</strong> functions φ j(k) such that the Fourier<br />

coefficients are close to the Fourier coefficients <strong>of</strong> the Eisenstein series E(z,s)<br />

where λ j(k) = s(1 − s).<br />

The typical difference between the Fourier coefficients at weight k and<br />

weight 0 are in both cases basically <strong>of</strong> order k; for the forms in C(k), the<br />

distance between λ j(0) and the corresponding discrete eigenvalue at weight 0<br />

is much smaller than k.<br />

The “generic” in the statement above means that we exclude certain places<br />

where the families φ j(k) changes character between C(k) and E(k). In these<br />

cases we have a situation <strong>of</strong> almost multiplicity 2, and in too coarse resolution<br />

it actually looks like the analytic families intersect.<br />

Note that Weyl’s law for non-trivial η-multiplier and a fixed non zero weight<br />

91


k on PSL(2,Z) is<br />

T 2 T<br />

<br />

<br />

Nk(T ) = ♯{R ≤ T } = − ln1<br />

− e<br />

12 π<br />

kπ 6 i<br />

<br />

<br />

+ S(T ) + O(1),<br />

for T ≥ 1 (cf. [39, p. 466]). One would like to say that the last two terms are<br />

bounded uniformly in k, and indeed one can easily show that the O(1) term<br />

is uniformly bounded in k. It seems plausible that, using techniques from<br />

[41, 38, 39], one should be able to prove that S(T ) is also uniformly bounded<br />

in k as k → 0. Details in that direction might appear in a later paper, but<br />

for now we simply make an assumption (heuristic, but backed by numerical<br />

evidence) that S(T ) is uniformly bounded in k. We then have the following<br />

(conjecture):<br />

Nk(T ) = N0(T ) − T<br />

π ln<br />

<br />

<br />

1 − e kπ 6 i<br />

<br />

<br />

+ O(1) = N0(T ) − T<br />

π ln<br />

<br />

kπ<br />

+ O(1),<br />

6<br />

(2.22)<br />

as k → 0, where<br />

T 2<br />

N0(T ) =<br />

12 .<br />

We have computed Nk(T ) experimentally for k = 10− j , j = 4,...,12, and Figure<br />

2.3 shows a picture <strong>of</strong> the experimental values compared to the theoretical<br />

values obtained by (2.22) (with the O(1) term neglected) and the difference<br />

indeed seems to be constant over this range.<br />

From the form <strong>of</strong> the Weyl law above, we also see that the successive spac-<br />

π<br />

ings ∆n(k) = Rn+1(k)−Rn(k) should look about like<br />

|ln kπ 6 | , i.e. ∆n(k) ≈ π<br />

|ln kπ 6 | .<br />

Figure 2.4 provides a nice illustration <strong>of</strong> this fact, where it is clearly seen that<br />

the average spacings are almost constant for small k and this constant is pro-<br />

. The mean gap over 10 typical cases (at R ≈ 10) turned out<br />

portional to 1<br />

|lnk|<br />

to be 0.118 at k = 10−9 , 0.109 at k = 10−10 and 0.100 at k = 10−11 . The rel-<br />

ative quotients, ∆n(k) <br />

kπ ln , 6 are 2.52, 2.58, 2.60, respectively, and it is not<br />

inconceivable that one obtains 3.14 in the (logarithmic) limit.<br />

Level repulsion<br />

From figures like 2.4 it one may be tempted to think that there are horizontal<br />

lines corresponding to cusp forms at weight 0 which crosses the lines that are<br />

going down (i.e. corresponding to the Eisenstein series at weight 0). This is<br />

not the case! If we look closer we will see that there is actually “level repulsion”<br />

here, i.e. the horizontal “cusp-form-line” turns down at the crossing and<br />

becomes an “Eisenstein-series-line”, and the previous “Eisenstein-series-line”<br />

turns into a “cusp-form-line”. More precisely formulated: if there is a “near<br />

crossing” at the weight k0 close to the eigenvalue R0 ≈ R j(k0) ≈ R j+1(k0), then<br />

92


Figure 2.3: Plot <strong>of</strong> Weyl’s law with constant T = 20 and weight k → 0<br />

Nk(20)<br />

200<br />

180<br />

160<br />

140<br />

120<br />

100<br />

80<br />

60<br />

3 4 5 6 7 8 9 10 11 12<br />

−log10(k)<br />

Nk(20)<br />

Experimental data<br />

Figure 2.4: Section <strong>of</strong> eigenvalues with 9 ≤ R ≤ 14, and 1E − 9 ≤ k ≤ 1E − 7.<br />

14<br />

13<br />

12<br />

11<br />

10<br />

9<br />

0 1e-08 2e-08 3e-08 4e-08 5e-08 6e-08 7e-08 8e-08 9e-08 1e-07<br />

there are two analytic families φ j(k) and φ j+1(k) such that for some δ > ε > 0:<br />

<br />

φ j(k) ∈<br />

C(k),<br />

E(k),<br />

k ∈ [k0 + ε,k0 + δ],<br />

k ∈ [k0 − δ,k0 − ε],<br />

and<br />

φ j+1(k) ∈<br />

<br />

E(k), k ∈ [k0 + ε,k0 + δ],<br />

C(k), k ∈ [k0 − δ,k0 − ε],<br />

and in the interval (k0 − ε,k0 + ε) both families display a mixing between the<br />

two types E(k) and C(k). In fact, the Fourier coefficients <strong>of</strong> φ j+1 converge<br />

93


(as k → k0 − ε) to values close to the Fourier coefficients <strong>of</strong> φ j for k > k0 +<br />

ε and vice versa. Since the two functions also need to be orthogonal it is<br />

clear that the Fourier coefficients exhibit “wild” behaviour in the small interval<br />

surrounding the “near crossing”. Note also that, as k → 0, all φ j(k) ∈ E(k) and<br />

R j(k) → 0.<br />

See Table 2.9 for examples <strong>of</strong> Fourier coefficients corresponding to eigenfunctions<br />

<strong>of</strong> types E(k) and C(k), close to an avoided crossing at weight<br />

k=9.044605824E-8. Table 2.10 illustrates the agreement between the Fourier<br />

coefficients <strong>of</strong> a more generic cusp form in E(k) and the corresponding coefficients<br />

<strong>of</strong> the Eisenstein series (appropriately normalized) at weight 0. The<br />

level <strong>of</strong> agreement is striking to put it mildly; likewise in Table 2.9 for the<br />

C(k) eigenfunction. The “1 for 1” nature <strong>of</strong> this convergence in the presence<br />

<strong>of</strong> a limiting continuous spectrum seems not to have been suspected earlier.<br />

Cf. [41, thm. 6.6 and cor. 6.9] and [40].<br />

The fact that the system seems to avoid accidental degeneracies by means<br />

<strong>of</strong> level repulsion and avoided crossings is in agreement with the Wigner-von<br />

Neumann theorem, cf. [101].<br />

2.9.3 Lifts at weight 1<br />

As we saw in Section 2.4.7 we could prove the existence <strong>of</strong> certain Hecke<br />

relations at weight 1 (e.g. (2.16), (2.17) and 2.18)). Tables 2.2 and 2.3 contain<br />

numerical verifications <strong>of</strong> these relations. Table 2.4 contains a list <strong>of</strong> computed<br />

eigenvalues on M Γ0(1),v 2 η,1,λ , and the eigenvalues corresponding<br />

to cosine CM-forms are indicated. In these cases, we have computed the actual<br />

error since we know the exact eigenvalues:<br />

Rk =<br />

2πk<br />

<br />

ln 7 + 2 √ , k ∈ Z<br />

12<br />

+ .<br />

Note that the actual error is in general much smaller than the error-parameter<br />

which is basically H(Y1,Y2) = |c(2) − c ′ (2)| + |c(3) − c ′ (3)| + |c(4) − c ′ (4)|,<br />

where c(n) is computed with Y1 and c ′ (n) with Y2.<br />

2.9.4 Half integer weight<br />

We now consider the case <strong>of</strong> Γ0(4) and the θ-multiplier system. The aim <strong>of</strong><br />

our investigation in this case was to study the Shimura lift, and in particular to<br />

investigate the dimensions <strong>of</strong> the spaces <strong>of</strong> half integer weight forms.<br />

Our numerical results indicate that the following conjectures are true:<br />

• If M(Γ0(2),0,2R) is an even (cosine) oldspace, then M(Γ0(4),vθ , 1<br />

2 ,R) has<br />

dimension 2. (Cf. our earlier remark about this in §2.7.1.)<br />

94


Figure 2.5: Level repulsion at R = 13.779...<br />

R<br />

R<br />

R<br />

14.1<br />

14.0<br />

13.9<br />

13.8<br />

13.7<br />

13.6<br />

13.5<br />

1.0e-03 1.2e-03 1.4e-03 1.6e-03 1.8e-03<br />

13.7805<br />

13.7800<br />

13.7795<br />

13.779760<br />

13.779755<br />

13.779750<br />

(a)<br />

Weight k<br />

1.6360e-03 1.6366e-03 1.6372e-03 1.6378e-03<br />

(c)<br />

Weight k<br />

1.636794e-03 1.636800e-03 1.636806e-03<br />

(e)<br />

Weight k<br />

R<br />

R<br />

R<br />

13.820<br />

13.800<br />

13.780<br />

13.760<br />

13.77980<br />

13.77975<br />

13.77970<br />

13.7797515<br />

13.7797510<br />

13.7797505<br />

1.60e-03 1.63e-03 1.66e-03 1.69e-03<br />

(b)<br />

Weight k<br />

1.63670e-03 1.63676e-03 1.63682e-03 1.63688e-03<br />

(d)<br />

Weight k<br />

1.6367962e-03 1.6367970e-03 1.6367978e-03<br />

(f)<br />

Weight k<br />

• In such 2-dimensional space, we can choose one (basis) function from the<br />

space V + , i.e. it has c1(n) = 0 if n ≡ 2,3 mod 4 (cf. §2.7.1).<br />

• Suppose that f ∈ M(Γ0(2),0,2R) is an even (cosine) newform, which satisfy<br />

f − 1<br />

<br />

2z = ε f (z), where ε = ±1. Then M(Γ0(4),vθ , 1<br />

2 ,R) is one dimensional.<br />

Furthermore, if the inverse Shimura lift <strong>of</strong> f , φ ∈ M(Γ0(4),vθ , 1<br />

2 ,R)<br />

has Fourier coefficients c(n) (with respect to ∞) then there are two cases:<br />

• If ε = 1 we can normalize φ with c(n) = 1 and we also have c(n) = 0<br />

for n ≡ 5 mod 8 (cf. R = 4.4614... in Table 2.9.4).<br />

• If ε = −1 then c(1) = 0, and we can normalize φ with c(2) = 1 and<br />

we also have c(n 2 ) = 0 for all n and c(n) = 0 for n ≡ 1 mod 8 (cf. R =<br />

6.0464... in Table 2.9.4).<br />

95


Figure 2.6: Level repulsion at R = 9.533...<br />

R<br />

R<br />

R<br />

10.0<br />

9.8<br />

9.6<br />

9.4<br />

9.2<br />

9.0<br />

9.533700<br />

9.533695<br />

9.533690<br />

9.53369535<br />

9.53369530<br />

9.53369525<br />

9.53369520<br />

9.53369515<br />

1.20e-08 1.60e-08 2.00e-08 2.40e-08 2.80e-08<br />

(a)<br />

Weight k<br />

2.43549e-08 2.43551e-08 2.43553e-08 2.43555e-08<br />

(c)<br />

Weight k<br />

2.4355247e-08 2.4355252e-08 2.4355257e-08<br />

(e)<br />

Weight k<br />

R<br />

R<br />

R<br />

9.53378<br />

9.53374<br />

9.53370<br />

9.53366<br />

9.53362<br />

9.5336956<br />

9.5336954<br />

9.5336952<br />

9.5336950<br />

9.533695264<br />

9.533695262<br />

9.533695260<br />

9.533695258<br />

9.533695256<br />

2.4350e-08 2.4353e-08 2.4356e-08 2.4359e-08<br />

(b)<br />

Weight k<br />

2.435523e-08 2.435525e-08 2.435526e-08 2.435528e-08<br />

(d)<br />

Weight k<br />

2.43552545e-08 2.43552547e-08 2.43552549e-08<br />

(f)<br />

Weight k<br />

Tables 2.5 and 2.9.4 contain examples <strong>of</strong> Fourier coefficients at weight 1<br />

2 and<br />

Table 2.8 contains a comparison <strong>of</strong> Fourier coefficients computed both from<br />

forms on Γ0(2) via (2.20) and computed directly. A supplement <strong>of</strong> Fourier<br />

coefficients for the weight 0 forms is available in Table 2.7.<br />

96


Table 2.1: Eigenvalues for N = 1, v = vη<br />

k = 5.0<br />

R |c(−1)| a H(y1,y2)<br />

3.66240686691 2E+3 8E-09<br />

5.77698688079 6E+3 1E-09<br />

6.64285171609 1E+4 2E-09<br />

7.82634704661 7E+4 8E-07<br />

8.66620831839 8E+4 1E-08<br />

9.45156176778 4E+4 4E-09<br />

10.21802876612 9E+4 2E-08<br />

10.65897262925 2E+5 5E-09<br />

11.27526358329 2E+5 2E-08<br />

12.15792337439 5E+6 2E-06<br />

12.55403510011 3E+5 3E-09<br />

13.00123950671 4E+4 2E-09<br />

13.67542640619 8E+5 8E-09<br />

13.71353384347 4E+6 4E-07<br />

14.47039277248 6E+5 1E-09<br />

15.03845367363 1E+6 6E-09<br />

15.39856858318 1E+6 1E-09<br />

15.85705128856 7E+4 7E-09<br />

16.14536205683 5E+6 2E-07<br />

16.45061260141 2E+6 1E-08<br />

16.93043847901 2E+7 2E-08<br />

17.51562192888 2E+5 4E-09<br />

17.59022138300 1E+6 6E-10<br />

18.13826107361 7E+6 2E-08<br />

18.32637702289 5E+5 5E-09<br />

18.76341585136 2E+6 5E-09<br />

19.16629116326 4E+6 8E-10<br />

19.67214438521 3E+7 2E-07<br />

19.68520099819 1E+6 3E-10<br />

20.00524829746 2E+6 4E-09<br />

20.38266630653 3E+6 3E-10<br />

20.67297062056 6E+6 4E-08<br />

20.97339376061 6E+6 8E-10<br />

k = 5.25<br />

a The normalization we have used here is the usual c(1) = 1.<br />

R |c(−1)| a H(y1,y2)<br />

3.68037312372 3E+3 3E-08<br />

5.82067054942 9E+3 6E-09<br />

6.63460520751 2E+4 3E-09<br />

7.90867228426 2E+5 9E-09<br />

8.61646891946 1E+5 8E-09<br />

9.56930344151 7E+4 5E-09<br />

10.15656706121 2E+5 2E-09<br />

10.70911890024 2E+5 4E-10<br />

11.34046324165 4E+5 2E-08<br />

12.11839521329 4E+6 2E-06<br />

12.65021958486 4E+5 9E-09<br />

13.02622821839 8E+4 2E-09<br />

13.56022943627 1E+6 1E-07<br />

13.87057635696 7E+5 4E-07<br />

14.48204838116 1E+6 7E-05<br />

15.09966704087 4E+7 1E-07<br />

15.38981845044 1E+6 4E-09<br />

15.94059443942 1E+4 3E-09<br />

16.09999759486 8E+6 7E-06<br />

16.52671557073 5E+6 1E-09<br />

16.90856097808 2E+7 5E-08<br />

17.53730159778 1E+6 3E-10<br />

17.74142373355 7E+6 4E-09<br />

18.02022951826 6E+6 9E-10<br />

18.37970066644 8E+5 2E-09<br />

18.90587158951 1E+7 8E-09<br />

19.09726131554 1E+7 5E-10<br />

19.66894569593 5E+6 3E-10<br />

19.73195996101 7E+6 4E-10<br />

20.12609436572 2E+6 3E-10<br />

20.35571778301 1E+7 7E-10<br />

20.71020380483 6E+6 8E-09<br />

20.88321504381 1E+7 2E-09<br />

97


Table 2.2: Fourier coefficients for a CM-form<br />

f ∈ M(Γ0(1),v 2 η,1,4.770984191561)<br />

n c(n)/c(0) a Error<br />

0 1.755930576575<br />

1 1.000000000000<br />

2 −1.755930576574 c(27) 0.4E-08<br />

3 1.571810322167 c(40) 0.1E-08<br />

4 1.755930576575 c(53) 0.7E-08<br />

5 −1.770268323978 c(66) 0.3E-09<br />

6 −2.474798320759 c(79) 0.1E-08<br />

7 0.000000000000 c(92) 0.4E-08<br />

8 0.346240855507 c(105) 0.5E-08<br />

9 3.510179255561 c(118) 0.3E-08<br />

10 1.116593241680 c(131) 0.4E-08<br />

11 −0.000000000001 c(144) 0.3E-08<br />

12 0.000000000001 c(157) 0.1E-07<br />

13 −3.019229958496 c(170) 0.8E-08<br />

14 −1.186431979458 c(183) + c(1) 0.3E-08<br />

15 −3.079783541463 c(196) 0.5E-08<br />

−c(n)c(−1)/R 2 /c(0)<br />

-1 5.055064268188 b 0.1E-11<br />

-2 −11.180067729976 c(21) 0.4E-07<br />

-3 0.000000000001 c(32) 0.2E-07<br />

-4 16.472675660354 c(43) 0.2E-08<br />

-5 16.729030199659 c(54) 0.5E-07<br />

-6 13.098490835617 c(65) 0.5E-08<br />

-7 16.046414105740 c(76) 0.1E-07<br />

-8 −0.000000000023 c(87) 0.1E-07<br />

-9 13.340740291248 c(98) 0.3E-07<br />

-10 0.000000000006 c(109) 0.4E-07<br />

-11 2.151215034503 c(120) 0.7E-08<br />

-12 2.878852009050 c(131) − c(1) 0.2E-07<br />

-13 0.000000000039 c(142) 0.6E-08<br />

-14 7.496795049955 c(153) 0.4E-07<br />

-15 −12.830881007618 c(164) 0.4E-07<br />

a This quotient is deduced from formula (2.14) or (2.15) on p.<br />

79.<br />

b c(−1) 2 = −R2c(0)(c(10) − c(0))<br />

98


Table 2.3: Fourier coefficients for a non-CM-form<br />

f ∈ M(Γ0(1),v 2 η,1,3.66240686698667)<br />

n c(n)/c(0) a Error<br />

0 −1.352193685534<br />

1 1.000000000000<br />

2 −1.697113317091 c(27) 0.5E-08<br />

3 −0.057989599353 c(40) 0.4E-07<br />

4 2.461764397786 c(53) 0.1E-07<br />

5 0.510856433057 c(66) 0.6E-08<br />

6 −0.952325903762 c(79) 0.8E-08<br />

7 −1.660630683908 c(92) 0.2E-07<br />

8 −2.343382246022 c(105) 0.9E-08<br />

9 1.271097206907 c(118) 0.1E-07<br />

10 −0.203512820511 c(131) 0.6E-08<br />

11 2.110622602834 c(144) 0.5E-08<br />

12 2.170616908700 c(157) 0.4E-08<br />

13 0.449799127363 c(170) 0.4E-07<br />

14 0.612654661780 c(183) + c(1) 0.4E-07<br />

15 −1.684453740441 c(196) 0.1E-07<br />

16 0.400312170289 c(209) 0.2E-07<br />

17 −2.868395110060 c(222) 0.1E-07<br />

18 −1.931595991172 c(235) 0.2E-07<br />

19 −0.591212766919 c(248) 0.1E-07<br />

20 0.792151138999 c(261) 0.9E-08<br />

21 −1.717242193922 c(274) 0.3E-07<br />

22 1.369169138277 c(287) 0.2E-07<br />

23 2.007854712832 c(300) 0.8E-09<br />

24 0.447826147902 c(313) 0.1E-07<br />

25 3.051006373828 c(326) 0.4E-08<br />

26 −0.032419986064 c(339) 0.6E-08<br />

27 1.255081531820 c(352) + a(2) 0.2E-07<br />

28 −0.707047087424 c(365) 0.3E-07<br />

29 1.272283355260 c(378) 0.1E-07<br />

30 −0.184187214400 c(391) 0.2E-07<br />

a This quotient is deduced from formula (2.14) or<br />

(2.15) on p. 79.<br />

99


Table 2.4: Eigenvalues for M(PSL(2,Z),1,v 2 η)<br />

R H(y1,y2) True error a<br />

2.38549209578045 1E-12 1E-15 b<br />

3.66240686698667 1E-11<br />

4.77098419156091 1E-13 1E-14 b<br />

5.77698688078694 6E-12<br />

6.64285171613711 1E-11<br />

7.15647628734173 1E-11 4E-13 b<br />

7.82634704540775 1E-13<br />

8.66620831896793 7E-14<br />

9.45156176783224 7E-12<br />

9.54196838312186 2E-13 6E-14 b<br />

10.21802876776059 3E-13<br />

10.65897262920241 2E-12<br />

11.27526358349387 2E-13<br />

11.92746047890219 7E-11 5E-14 b<br />

12.15792337422149 7E-13<br />

12.55403509998720 1E-13<br />

13.00123950642372 9E-13<br />

13.67542640643589 2E-13<br />

13.71353384358095 1E-12<br />

14.31295257468268 1E-12 1E-14 b<br />

14.47039277253940 1E-12<br />

15.03845367358721 2E-13<br />

15.39856858348441 1E-12<br />

15.85705128717333 3E-10<br />

16.14536205734475 3E-11<br />

16.45061260131967 4E-11<br />

16.69844467046304 4E-12 1E-13 a<br />

16.93043847896222 4E-15<br />

17.51562192885174 2E-10<br />

17.59022138305996 1E-10<br />

18.13826107340244 3E-11<br />

18.32637702205910 4E-11<br />

18.76341585146817 1E-11<br />

a For CM-forms, the true error is computed with<br />

respect to the eigenvalue Rk = 2πk<br />

ln(η0) , where<br />

η0 = 7 + 2 √ 12.<br />

b These forms correspond to CM-forms.<br />

100


Table 2.5: Fourier coefficients for<br />

f1,2 ∈ M(Γ0(4), 1<br />

2 ,6.889875675945)<br />

Fourier coefficients for f1 ∈ V + , observe that a(n) = 0 for n ≡ 2,3 mod 4<br />

n a(n)<br />

4 0.84219769675471<br />

5 0.18355821406443<br />

8 0.56907998524429<br />

9 −0.33045049673565<br />

12 −0.41296169213831<br />

13 0.60153537988057<br />

16 0.30482066289397<br />

17 −0.88689598690620<br />

20 0.41418282092059<br />

21 0.50212917023175<br />

23 −0.00000000000110<br />

24 −1.04429548341249<br />

25 0.28984678984391<br />

n a(−n)<br />

|c(4)c(9) − c(36)| = 0.2E − 08<br />

3 1.01825299171456<br />

4 2.18968040385979<br />

7 2.06305218095270<br />

8 −1.17157116610978<br />

11 −1.02718719694121<br />

12 2.29759751514531<br />

15 −4.08935474990375<br />

16 3.39248165496032<br />

19 −1.31804633824673<br />

20 1.87725570455517<br />

23 −2.22265197818114<br />

24 0.58816620330140<br />

Fourier coefficients for f2 ∈ V + , observe that a(1) = 0, and we set c(2) = 1<br />

n a(n) n a(n)<br />

0 0.000000000000 19 −0.936283350934<br />

1 0.000000000000 20 0.192087703124<br />

2 1.000000000000 21 1.247834928335<br />

3 −0.725665465042 22 0.411443268212<br />

4 −0.710754741008 23 −0.018564690418<br />

5 0.456158224759 24 −1.297582830232<br />

6 −1.835059236817 25 0.000000000000<br />

7 0.481289183972 26 −0.610656711780<br />

8 0.707106781187 27 −0.179166638197<br />

9 0.000000000000 28 0.340322845698<br />

10 −0.651049038069 29 0.001563849679<br />

11 −0.470914040036 30 1.462745418892<br />

12 −0.513122971204 31 0.400763051265<br />

13 1.494868058150 32 0.095523702475<br />

14 0.968484734380 33 0.000000000000<br />

15 −0.945563574521 34 −0.879782299796<br />

16 −1.101175502961 35 −1.833603623066<br />

17 0.000000000000 36 0.234869257223<br />

18 0.824250041644<br />

101


Table 2.6: Fourier coefficients for f ∈ M Γ0(4), 1<br />

2 ,R<br />

102<br />

R = 4.461438243496<br />

n a(n)<br />

0 0.00000000000000<br />

1 1.00000000000000<br />

2 0.63334968449036<br />

3 0.63517832947402<br />

4 −0.70710678118667<br />

5 −0.00000000000003<br />

6 1.28035706400142<br />

7 −0.90756258916698<br />

8 −0.44784585676555<br />

9 0.52643872643776<br />

10 −0.57763498966972<br />

11 1.12485377915641<br />

12 −0.44913890403389<br />

13 0.00000000000001<br />

14 0.48078071833327<br />

15 −1.58539012005784<br />

16 0.50000000000015<br />

17 −0.14882069214483<br />

18 1.06474902295605<br />

19 0.21268916863632<br />

20 0.00000000000003<br />

21 0.00000000000002<br />

22 −1.56248056455209<br />

23 0.85478501960318<br />

24 −0.90534916229569<br />

25 0.45696099733973<br />

36 −0.37224839333969<br />

|c(4)c(9) − c(36)| = 4E − 12<br />

R = 6.046497437542<br />

n a(n)<br />

0 0.000000000000<br />

1 0.000000000000<br />

2 1.000000000000<br />

3 1.770795863371<br />

4 0.000000000029<br />

5 1.331470494003<br />

6 −0.749395507674<br />

7 0.074369313625<br />

8 0.707106781209<br />

9 −0.000000000016<br />

10 −1.214451367832<br />

11 0.620756243833<br />

12 1.252141763204<br />

13 −1.123686021586<br />

14 −0.711286031630<br />

15 0.964706032413<br />

16 −0.000000000112<br />

17 −0.000000000108<br />

18 −0.128644484609<br />

19 −0.499664871594<br />

20 0.941491814843<br />

21 −0.447011950327<br />

22 −0.805501416928<br />

23 1.236716356721<br />

24 −0.529902643583<br />

25 0.000000008756<br />

36 −0.000000000168<br />

288 −0.064322242377<br />

|c(2 · 3 2 )c(2 · 4 2 ) − c(2 · 12 2 )| = 8E − 11


Table 2.7: Supplemental table <strong>of</strong> Fourier coefficients for M(Γ0(2),0,R)<br />

R 8.92287648699174<br />

n A(n)<br />

2 −0.70710678118654<br />

3 1.10378899562734<br />

4 0.49999999999993<br />

5 0.90417459283958<br />

6 −0.78049668380711<br />

7 0.82934246755499<br />

8 −0.35355339059330<br />

9 0.21835014686776<br />

10 −0.63934798597339<br />

R 13.77975135189073<br />

(even wrt z ↦→ − 1 2z )<br />

n A(n)<br />

2 2.96351804031448<br />

3 0.24689977245401<br />

4 3.59139177031902<br />

5 0.73706038534834<br />

6 0.73169192981688<br />

7 −0.26142007576500<br />

8 2.60064131148226<br />

9 −0.93904050235826<br />

10 2.18429174878080<br />

R 12.0929948750786<br />

n A(n)<br />

2 0.70710678118655<br />

3 −0.70599475399569<br />

4 0.49999999999999<br />

5 −0.79974825694039<br />

6 −0.49921367803249<br />

7 −1.71337067862845<br />

8 0.35355339059328<br />

9 −0.50157140733054<br />

10 −0.56550741572467<br />

R 13.77975135189073<br />

(odd wrt z ↦→ − 1 2z )<br />

n A(n)<br />

2 0.13509091556820<br />

3 0.24689977245398<br />

4 −0.79070303958101<br />

5 0.73706038534830<br />

6 0.03335391631439<br />

7 −0.26142007576538<br />

8 −1.36013067551284<br />

9 −0.93904050236089<br />

10 0.09957016228575<br />

103


Table 2.8: Comparison <strong>of</strong> Fourier coefficients A(n) computed directly for<br />

M(Γ0(N),0,R) (N = 1,2), vs Â(n) computed on M Γ0(4), 1 1<br />

2 , 2R and using<br />

(2.20).<br />

n Â(n) A(n) |A(n) − Â(n)|<br />

f ∈ M(Γ0(2),0,8.922876486992) (t = 1)<br />

2 −0.70710678118665 −0.707106781186 0.6E-12<br />

3 1.10378899562739 1.103788995627 0.4E-12<br />

5 0.90417459283969 0.904174592840 0.3E-12<br />

f ⊆ M(Γ0(1),0,13.779751351891) (t = 1)<br />

2 1.54930447794126 1.549304477941 0.7E-12<br />

3 0.24689977245398 0.246899772454 0.3E-12<br />

4 1.40034436536892 1.400344365369 0.2E-12<br />

5 0.73706038534387 0.737060385348 0.4E-11<br />

6 0.38252292109716 0.382522923066 0.2E-08<br />

f ∈ M(Γ0(1),0,13.779751351891) (t = 2)<br />

3 0.24689977245437 0.246899772454 0.8E-13<br />

5 0.73706038535004 0.737060385348 0.2E-11<br />

7 −0.26142007624377 −0.261420075765 0.5E-09<br />

9 −0.93904050238904 −0.939040502362 0.3E-10<br />

f ∈ M(Γ0(2),0,12.092994875079) (t = 2)<br />

3 −0.70599475379863 −0.705994753996 0.2E-09<br />

5 −0.79974825694696 −0.799748256940 0.7E-11<br />

7 −1.71337067860377 −1.713370678628 0.2E-10<br />

9 −0.50157140750090 −0.501571407330 0.2E-09<br />

In segment 3, the calculation is based on (2.20) and the second portion <strong>of</strong> Table 2.5.<br />

104


Table 2.9: Comparison <strong>of</strong> Fourier coefficients for weights k =9.044605824E-08 and<br />

k = 0 near an “avoided crossing”.<br />

Corresponds to the cusp form<br />

k<br />

R<br />

c(2)<br />

c(3)<br />

c(4)<br />

c(5)<br />

c(6)<br />

9.0446058240E-08<br />

13.77975135189074<br />

1.54930480559976<br />

0.24689988546553<br />

1.40034433555250<br />

0.73706067260516<br />

0.38252272069428<br />

Corresponds to the Eisenstein series<br />

k<br />

R<br />

c(2)<br />

c(3)<br />

c(4)<br />

c(5)<br />

c(6)<br />

9.0446058240E-08<br />

13.77975135138225<br />

−2.06525760334129<br />

−1.72891679536648<br />

2.97153986917404<br />

−2.02747287754385<br />

3.41008729668221<br />

0<br />

13.77975135189074<br />

1.54930447794069<br />

0.24689977245411<br />

1.40034436536841<br />

0.73706038534787<br />

0.38252292306557<br />

0<br />

13.77975135138225<br />

−1.98398933080188<br />

−1.68449330640991<br />

2.93621366473571<br />

−1.96531634618530<br />

3.34201674772446<br />

Difference<br />

0.3E-06<br />

0.1E-06<br />

0.1E-06<br />

0.2E-06<br />

0.2E-06<br />

Difference<br />

0.8E-01<br />

0.4E-01<br />

0.4E-01<br />

0.6E-01<br />

0.7E-01<br />

The coefficients for the Eisenstein series at k = 0 were computed using (2.19), i.e.:<br />

c(2) = 2cos(Rln2)<br />

c(3) = 2cos(Rln3)<br />

c(4) = 1 + 2cos(Rln4)<br />

c(5) = 2cos(Rln5)<br />

c(6) = 2cos(Rln6) + 2cos(R(ln3 − ln2))<br />

Table 2.10: Comparison <strong>of</strong> Fourier coefficients for weights k =9.044605824E-08 and<br />

k = 0 “far” from an ”avoided crossing”. The weight 0 coefficients were computed<br />

using the formulas in Table 2.9.<br />

k 9.0446058240E-08 0<br />

Difference<br />

R<br />

c(2)<br />

c(3)<br />

c(4)<br />

c(5)<br />

c(6)<br />

13.62696884857618<br />

−1.99957085683552<br />

−1.48069687587703<br />

2.99828354611637<br />

−1.99647405201235<br />

2.96075820067617<br />

13.62696884857618<br />

−1.99957081810438<br />

−1.48069680342062<br />

2.99828345661464<br />

−1.99647406885962<br />

2.96075811858031<br />

0.4E-07<br />

0.7E-07<br />

0.7E-07<br />

0.2E-07<br />

0.8E-07<br />

105


3 Computation <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong> on<br />

Non-congruence Subgroups <strong>of</strong> the Modular<br />

Group<br />

3.1 Introduction<br />

Previously, in Chapter 1 we considered the computation <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong><br />

on certain subgroups <strong>of</strong> PSL(2,Z), the so called Hecke congruence subgroups,<br />

and in Chapter 2 we extended this work to the case <strong>of</strong> arbitrary real weight<br />

and non-trivial multiplier system. In this chapter we will see that the methods<br />

detailed in the previous chapters also work for non-congruence subgroups.<br />

There is a general belief that there exist no (or very few) <strong>Maass</strong> <strong>waveforms</strong><br />

on a generic Fuchsian group. This is indicated by the results <strong>of</strong> Phillips and<br />

Sarnak, cf. [80, 81]. Recent results (cf. [30]) indicate that there are continuous<br />

families <strong>of</strong> groups, on which there exists (at least) one <strong>Maass</strong> waveform.<br />

These families seem to be “arbitrary” in the sense that they are not arithmetic<br />

and possess no obvious symmetries.<br />

If a group Γ has <strong>Maass</strong> <strong>waveforms</strong>, then these descend to all subgroups <strong>of</strong><br />

Γ. As a <strong>Maass</strong> waveform on the subgroup we will call it an oldform. We will<br />

use the term newform to denote the <strong>Maass</strong> <strong>waveforms</strong> on a group G which<br />

are not from a supergroup G1 (G G1 ⊆ PSL(2,Z)). This is an analog <strong>of</strong><br />

the well-known concepts <strong>of</strong> oldforms and newforms for congruence groups<br />

(cf. [5]).<br />

The groups which we will study here are all subgroups <strong>of</strong> the modular<br />

group, PSL(2,Z), and hence arithmetic, but even in this case the existence<br />

or nonexistence <strong>of</strong> newforms is in general not known.<br />

The existence <strong>of</strong> an infinite number <strong>of</strong> newforms is known on congruence<br />

subgroups and cycloidal subgroups (having only one cusp), and other subgroups<br />

which possess certain symmetries (e.g. a reflectional symmetry which<br />

preserves cusp classes). The results presented in this thesis indicate that there<br />

are no other examples <strong>of</strong> subgroups <strong>of</strong> index ≤ 11 which have newforms.<br />

We will give a general method that produces all subgroups with a given<br />

index in PSL(2,Z), and then we will see how to adapt the previous algorithm<br />

for computation <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong> (cf. Section 1.3) to such a group, and<br />

finally we will provide examples <strong>of</strong> the experimental results.<br />

The author would also like to express gratitude to Andreas Strömbergsson<br />

107


for very helpful discussions, especially regarding Proposition 3.2.1.<br />

3.2 Subgroups <strong>of</strong> the Modular Group<br />

Remember that the modular group, PSL(2,Z), is defined as<br />

<br />

a b ad<br />

PSL(2,Z) =<br />

− bc = 1,a,b,c,d ∈ Z ,<br />

c d<br />

with a (closed) fundamental domain<br />

<br />

F1 = z = x + iy ∈ H <br />

<br />

1<br />

|z| ≥ 1, |x| ≤ .<br />

2<br />

For a given positive integer N the principal congruence subgroup, Γ(N), is<br />

defined by<br />

<br />

a b<br />

Γ(N) =<br />

∈ PSL(2,Z)<br />

c d<br />

<br />

c ≡ b ≡ 0 mod N,a ≡ d ≡ 1 mod N ,<br />

and any subgroup <strong>of</strong> PSL(2,Z) containing some Γ(N) is called a congruence<br />

subgroup. If G ⊆ PSL(2,Z), and N is the smallest integer such that Γ(N) ⊆ G<br />

then we say that N is the level <strong>of</strong> G. A particularly interesting and well understood<br />

class <strong>of</strong> congruence subgroups are the Hecke congruence subgroups,<br />

Γ0(N),<br />

<br />

a b<br />

Γ0(N) =<br />

∈ PSL(2,Z)<br />

c d<br />

<br />

c ≡ 0 mod N .<br />

Presentations, constructions, modular forms and spectral theory are more or<br />

less well known and understood topics on congruence groups in general and<br />

on Hecke congruence groups in particular, cf. e.g. [83, 24, 76, 84, 89, 51, 5,<br />

62, 12].<br />

Of the other types <strong>of</strong> subgroups, the cycloidal groups have been studied<br />

for a long time by Petersson, Millington, etc. (cf., e.g. [78, 79, 66, 67]), and<br />

other specific types <strong>of</strong> subgroups have also been studied by various people<br />

(cf., e.g. [22, 23, 26]). The subject <strong>of</strong> interest to these investigations however<br />

has mostly been in presenting and constructing the groups. The present<br />

knowledge <strong>of</strong> the theory <strong>of</strong> modular forms and spectral theory on these groups<br />

is very limited (cf., e.g. [7, 105, 106, 100]).<br />

It is a well-known fact that PSL(2,Z) is generated by the elements S : z ↦→<br />

z + 1, and E : z ↦→ − 1<br />

z , and that the only relations are E2 = R 3 = Id, where<br />

108


R = ES. In terms <strong>of</strong> matrices we have<br />

<br />

S =<br />

1<br />

0<br />

1<br />

1<br />

, E =<br />

0<br />

1<br />

−1<br />

0<br />

, and R =<br />

0<br />

1<br />

−1<br />

1<br />

.<br />

Let Γ ⊆ PSL(2,Z) be a subgroup <strong>of</strong> finite index µ with closed fundamental<br />

domain F in H. Suppose that F has e2 and e3 inequivalent elliptic vertices <strong>of</strong><br />

order 2 and 3 respectively and κ inequivalent parabolic vertices (cusps). The<br />

genus g <strong>of</strong> the quotient orbifold, Γ\H, is given by<br />

g = 1 + 1<br />

2<br />

<br />

µ 1<br />

− κ −<br />

6 2 e2 − 2<br />

3 e3<br />

<br />

, (3.1)<br />

(cf. [67, (1.1)]) and we say that the type <strong>of</strong> Γ is (µ; g,κ,e2,e3). If µ j, j =<br />

1,...,κ are the widths <strong>of</strong> the cusps <strong>of</strong> F then µ = µ1 +···+ µκ. (This formula<br />

is easily verified by viewing F as µ copies <strong>of</strong> F1, the standard fundamental<br />

domain for PSL(2,Z),since then µ j is exactly the number <strong>of</strong> images <strong>of</strong> F1<br />

whose vertex at infinity belongs to the cusp class number j.) The generalized<br />

level <strong>of</strong> Γ is defined as the least common multiple <strong>of</strong> µ j<br />

κ j=1<br />

. It is known that<br />

if Γ is a congruence subgroup then the generalized level is equal to the level<br />

(cf. [104]).<br />

Theorem 3.1. (Millington [67])<br />

For every set <strong>of</strong> integers µ > 0, κ > 0, e2 ≥ 0, e3 ≥ 0, g ≥ 0 consistent with<br />

(3.1) there exists a subgroup <strong>of</strong> the modular group <strong>of</strong> type<br />

(µ; g,κ,e2,e3).<br />

3.2.1 Permutations and subgroups <strong>of</strong> PSL(2,Z)<br />

There is a convenient way to represent any subgroup G <strong>of</strong> PSL(2,Z) <strong>of</strong> finite<br />

index µ in terms <strong>of</strong> a pair <strong>of</strong> permutations <strong>of</strong> a set <strong>of</strong> µ elements. Geometrically<br />

these permutations can be viewed as describing the side pair identifications<br />

in a fundamental domain for G formed by µ distinct PSL(2,Z)-translates<br />

<strong>of</strong> F1 (cf. Remark 3.2.1 for a precise description <strong>of</strong> this).<br />

For a positive integer µ we use Sµ to denote the group <strong>of</strong> permutations <strong>of</strong><br />

µ-letters (e.g. {1,2,..., µ}). The notation here is that the action <strong>of</strong> a permutation<br />

π ∈ Sµ is from the left, i.e. for π,σ ∈ Sµ we write πσ( j) = π (σ( j)).<br />

A k−cycle is a permutation π ∈ Sµ with the property that there exist k distinct<br />

elements x1,x2,...,xk ∈ {1,2,..., µ} such that π(xi) = xi+1, i = 1,...,k−<br />

1 and π(xk) = x1. We will use the following notation: π = (x1x2 ···xk). Note<br />

that with this notation, (x1x2 ···xk) = (x2x3 ···xkx1) = ... .<br />

A general permutation π is completely described by either listing its values<br />

109


on all elements, i.e. {π( j)} µ<br />

j=1 or by listing all cycles, i.e.<br />

with the sets x j1,xj2,...,xjk j<br />

π = (x11 ···x1k1 )(x21 ···x2k2 )···(xn1 ···xnkn ),<br />

( j = 1,2,...,n) pairwise disjoint.<br />

Definition 3.2.1. If σE,σR ∈ Sµ we say that the pair (σE,σR) is legitimate if<br />

• σ 2 E = σ 3 R = 1, and<br />

• the group generated by σE and σR, 〈σE,σR〉, is transitive on {1,2,..., µ}.<br />

Two legitimate pairs (σE,σR) and (σ ′ E ,σ′ R ) are said to be equivalent, denoted<br />

by (σE,σR) ∼ (σ ′ E ,σ′ R ), if there exists a permutation π ∈ Sµ such that<br />

π −1 σEπ = σ ′ E, and π −1 σRπ = σ ′ R.<br />

If we also have π(1) = 1, we say that the pairs are equivalent (mod 1), which<br />

we denote by (σE,σR) ∼1 (σ ′ E ,σ′ R ).<br />

We have the following theorem.<br />

Theorem 3.2. (Millington [67, thm. 2])<br />

There is a one-to-one correspondence between equivalence classes (mod 1)<br />

(i.e. with respect to ∼1) <strong>of</strong> legitimate pairs (σE,σR) in Sµ and subgroups G<br />

<strong>of</strong> index µ in PSL(2,Z). Furthermore, G has type (µ; g,κ,e2,e3) and cusp<br />

widths µ j<br />

κ j=1<br />

if and only if<br />

• σE has e2 1−cycles (fixed elements),<br />

• σR has e3 1−cycles (fixed elements), and<br />

• σS = σRσE has κ disjoint cycles <strong>of</strong> lengths µ j, 1 ≤ j ≤ κ<br />

(and g is given by (3.1)).<br />

Remark 3.2.1. The subgroup, G, <strong>of</strong> PSL(2,Z) corresponding to the legitimate<br />

pair (σE,σR) is defined by<br />

G = {γ ∈ PSL(2,Z)|h(γ)(1) = 1},<br />

where h : PSL(2,Z) → Sµ is the anti-homomorphism (i.e. h(AB) = h(B)h(A),<br />

for all A,B ∈ PSL(2,Z) and h(Id) = Id) defined by h(E) = σE and h(R) = σR.<br />

In the implementation <strong>of</strong> the general pullback algorithm (cf. Section 1.3.2) we<br />

use the fact that a set <strong>of</strong> right coset representatives <strong>of</strong> G in PSL(2,Z) is given<br />

by µ<br />

Vj , where<br />

j=1<br />

h(Vj)(1) = j, j = 1,..., µ.<br />

Hence a fundamental domain for G\H is given by F = ⊔ µ<br />

j=1Vj (F1). One<br />

<br />

checks from the definitions that, for all j, h<br />

(1) = 1 so that for all<br />

<br />

VjSV −1<br />

h(S)( j)<br />

j, VjSV −1<br />

h(S)( j) ∈ G. This means that VjS(F1) can be mapped to Vh(S)( j) (F1) by<br />

a map in G. Let b be the geodesic ei π 3 + i[0,∞) so that b and S−1 (b) are two<br />

110


<strong>of</strong> the three sides <strong>of</strong> the geodesic polygon F1. Since Vj (F1) ∩VjS(F1) = Vj(b)<br />

it follows that there is an identification map in G which glues the side Vj(b)<br />

<strong>of</strong> Vj (F1) to the side V h(S)( j)S −1 (b) <strong>of</strong> V h(S)( j) (F1). Similarly, if b ′ denotes<br />

the side <strong>of</strong> the polygon F1 which contains the point i, then the side Vj(b ′ ) <strong>of</strong><br />

Vj (F1) is glued to the side V h(E)( j)(b ′ ) (provided with the “opposite direction”)<br />

<strong>of</strong> V h(E)( j) (F1).<br />

For the pullback procedure, we need to find the coset <strong>of</strong> the map T −1 , where<br />

T gives the pullback into F1, i.e. we need to find h(T −1 )(1). To do that, as<br />

we go through with the “flip-flop” through the generators E and S, in each<br />

step we both multiply the appropriate map (from the left) and the inverse <strong>of</strong> the<br />

corresponding permutation (from the left). When we are done, we have both T<br />

and h(T −1 ) and can thus determine the coset <strong>of</strong> T −1 . (The G-pullback is then<br />

given by the map V h(T −1 )(1)T , cf. Chapter 1, p. 43.)<br />

(In her thesis, Millington appears to get a connected polygon with identified<br />

sides by linking the permutations with the Poincaré fundamental polygon theorem<br />

in an appropriate way. See [66, pro<strong>of</strong> <strong>of</strong> thm. 1] for the cycloidal case.<br />

In our approach, we do not need a connected polygon and simply proceed<br />

algebraically by way <strong>of</strong> [67, thm. 1].)<br />

Theorem 3.3. Let (σE,σR) correspond to G and (σ ′ E ,σ′ R ) to G′ . Then G and<br />

G ′ are conjugate in PSL(2,Z) iff (σE,σR) ∼ (σ ′ E ,σ′ R ).<br />

Pro<strong>of</strong>. Let h and h ′ denote the anti-homomorphisms corresponding to (σE,σR)<br />

and (σ ′ E ,σ′ R ). Suppose that A ∈ SL(2,Z). We then know that<br />

AGA −1 = A γ ∈ PSL(2,Z) h(γ)(1) = 1 A −1<br />

= AγA −1 ∈ PSL(2,Z) h(γ)(1) = 1 <br />

= γ ∈ PSL(2,Z) h(A)h(γ)h(A) −1 (1) = 1 <br />

= γ ∈ PSL(2,Z) h ′′ (γ)(1) = 1 ,<br />

where h ′′ (γ) = h(A)h(γ)h(A) −1 ∈ Sµ corresponds to the legitimate pair (σ ′′<br />

E ,σ′′ R )<br />

(here σ ′′<br />

E = h(A)σEh(A) −1 and σ ′′<br />

R = h(A)σRh(A) −1 ). Suppose that G ′ = AGA−1 .<br />

This is equivalent with<br />

<br />

γ ∈ PSL(2,Z) h ′′<br />

(γ)(1) = 1 = γ ∈ PSL(2,Z) h ′<br />

(γ)(1) = 1 ,<br />

), which<br />

means that there exists a π ∈ Sµ with π(1) = 1 such that π−1 (σ ′ E ,σ′ R )π =<br />

and by theorem 3.2 we must have then have (σ ′′<br />

E ,σ′′<br />

R ) ∼1 (σ ′ E ,σ′ R<br />

(σ ′′<br />

E ,σ′′<br />

R ) and if we set τ = h(A)−1 π −1 we have τ −1 (σE,σR)τ = (σ ′ E ,σ′ R )<br />

where now τ(1) = h(A) −1 (1). Hence (σE,σR) ∼ (σ ′ E ,σ′ R ) .<br />

Conversely, suppose that τ ∈ Sµ and τ−1 (σE,σR)τ = (σ ′ E ,σ′ R ). Then, for<br />

any A ∈ SL(2,Z) such that h(A) −1 (1) = τ(1), we have<br />

G ′ = AGA −1 .<br />

111


This is easily seen by comparing h ′ (γ) = τ −1 h(γ)τ and h ′′ (γ) = h(A)h(γ)h(A) −1<br />

for γ ∈ PSL(2,Z).<br />

Remark 3.2.2. If G ′ = AGA−1 , A ∈ SL(2,Z) then if j = h(A) −1 (1) and σ j =<br />

(1 j), the pair (σ ′ E ,σ′ R ) defined by<br />

′<br />

σ E,σ ′ −1<br />

R = σ j (σE,σR)σj<br />

is easily seen to correspond to the group G ′ . Hence we see that all conjugates<br />

<strong>of</strong> G are obtained by conjugating the corresponding pair (σE,σR) by σ j for<br />

j = 2,..., µ.<br />

Definition 3.2.2. If G is a subgroup <strong>of</strong> PSL(2,Z) then we use (G) to denote<br />

the set <strong>of</strong> PSL(2,Z)-conjugates <strong>of</strong> G, (G) = AGA −1 A ∈ PSL(2,Z) .<br />

It is clear now that if (σE,σR) is a legitimate pair corresponding to G, then<br />

we get all groups in (G) by conjugating the pair (σE,σR) with all σ j = (1 j),<br />

j = 2,..., µ.<br />

Remark 3.2.3. We know that conjugate groups have the same spectrum, hence<br />

we only need to investigate one G in the set (G).<br />

The reflection operator J : z ↦→ −z, which we also represent by the matrix<br />

<br />

J =<br />

−1<br />

0<br />

0<br />

1<br />

,<br />

acts on E and R as<br />

E ∗ = JEJ = E, and R ∗ = JRJ = ER −1 E = ER 2 E.<br />

Hence, if G is represented by the legitimate pair (σE,σR), the group G∗ = JGJ<br />

corresponds to the legitimate pair<br />

<br />

.<br />

(σ ∗ E,σ ∗ R) = σE,σEσ 2 RσE<br />

We also observe that G ∗ has the same type as G and the spectrum <strong>of</strong> G and<br />

G ∗ are equal. In fact, it is obvious that if f is a function on G\H, then f |J is a<br />

function on G ∗ \H.<br />

Remark 3.2.4. Since one can easily compute the pair corresponding to G ∗<br />

from the pair corresponding to G we do not list (G ∗ ) separately in Table 3.1.<br />

But for each conjugacy class (G) we list the conjugacy class <strong>of</strong> (G ∗ ). Clearly<br />

either (G ∗ ) = (G) or (G ∗ ) = (G). If (G ∗ ) = (G) we can in some cases use<br />

the reflection to prove existence <strong>of</strong> newforms (cf Proposition 3.2.1 and Remark<br />

3.2.1).<br />

112


We will now give a precise statement about a sufficient condition that newforms<br />

exist on a subgroup <strong>of</strong> PSL(2,Z). This proposition gives an explanation<br />

for all newforms on noncongruence subgroups which we have found so far.<br />

Proposition 3.2.1. Suppose that G corresponds to the legitimate pair (σE,σR),<br />

and that JGJ = A −1 GA for some A ∈ SL(2,Z). Let P = {[p1],...,[pκ]} be the<br />

set <strong>of</strong> G-equivalence classes <strong>of</strong> cusps <strong>of</strong> G. Set α = AJ and assume that α 2 ∈ G<br />

and that α preserves all cusp classes. Then L 2 (G\H) splits into an “α-even”<br />

and “α-odd” part, and the α-odd part is spanned by <strong>Maass</strong> <strong>waveforms</strong>.<br />

Pro<strong>of</strong>. Extending [94, I, §§1-2] to PGL(2,R) (cf. Chapter 1, p. 21) we see that<br />

α gives us a modular correspondence (with d = 1). Since the cusp classes are<br />

preserved by α we can use [94, lemma 2.6] (extended to PGL(2,R)) to show<br />

that there exist positive constants γ1,...,γκ such that<br />

E j (αz,s) ≡ γ s jEj (z,s), for j = 1,...,κ, (3.2)<br />

where E j (z,s) is the Eisenstein series for G\H associated with [p j]. Since<br />

α2 <br />

∈ G we have E j α2z,s ≡ E j (z,s), and hence, in fact γ1 = γ2 = ··· = γκ =<br />

1. L2 (G\H) splits into an “even” and an “odd” part by the action <strong>of</strong> α;<br />

L 2 (G\H) = L 2 + (G\H) ⊕ L 2 − (G\H),<br />

(an orthogonal sum, since the action <strong>of</strong> α is self adjoint, by an argument similar<br />

to [38, pp. 470-472] but easier) where<br />

L 2 ± (G\H) = f ∈ L 2 (G\H) f (αz) ≡ ± f (z) .<br />

It is known that L 2 (G\H) has a spectral decomposition as an orthogonal sum<br />

L 2 (G\H) = C ⊕D, where C (“the continuous part”) is spanned by the Eisenstein<br />

series, and D is spanned by <strong>Maass</strong> <strong>waveforms</strong>. By (3.2) (where γ j = 1)<br />

we have C ⊂ L 2 + (G\H). Hence L 2 − (G\H) ⊂ D, i.e. L 2 − (G\H) is spanned<br />

by <strong>Maass</strong> <strong>waveforms</strong>. It will be seen momentarily (cf. Remark 3.2.5) that, in<br />

many cases, the space L 2 − (G\H) can be seen to be infinite-dimensional.<br />

Remark 3.2.5. In each case which we have investigated and for which Prop. 3.2.1<br />

applies, the group G does not allow any nontrivial supergroups G ′ (G G ′ <br />

PSL(2,Z)). Hence, in these cases, Prop. 3.2.1 very easily implies the existence<br />

<strong>of</strong> newforms for G.<br />

Indeed, let V ⊂ H be a small disk such that γV ∩V = /0 for all γ ∈ GL(2,Z);<br />

write PSL(2,Z) = ⊔n j=1Gγ j, and let f0 be any function {1,...,n} → C with<br />

∑ n j=1 f0( j) = 0, f0 ≡ 0. Define f : H → C by<br />

⎧<br />

⎪⎨ f0( j), if z ∈ Gγ jV,<br />

f (z) = − f0( j), if z ∈ Gαγ<br />

⎪⎩<br />

jV,<br />

0, otherwise.<br />

113


Then f is G-invariant and f ∈ L 2 − (G\H) (using αG = Gα = GαG). Further-<br />

more, if f1 ∈ L2 (PSL(2,Z)\H) then<br />

<br />

<br />

f f1dµ =<br />

G\H<br />

=<br />

=<br />

=<br />

= 0.<br />

n <br />

∑ + f f1dµ<br />

j=1 γ jV αγ jV<br />

n <br />

n <br />

∑ f0( j) f1dµ − ∑ f0( j) f1dµ<br />

j=1 γ jV j=1 αγ jV<br />

n <br />

n <br />

∑ f0( j) f1dµ − ∑ f0( j) f1dµ<br />

j=1 V j=1 αV<br />

n<br />

f1dµ − f1dµ ∑ f0( j)<br />

V<br />

αV j=1<br />

(We used αγ jV = αγ jα −1 αV , where αγ jα −1 ∈ PSL(2,Z)). Hence f is orthogonal<br />

to the full oldform space (by our assumption on G). Hence, expanding f<br />

as an L 2 -sum <strong>of</strong> <strong>Maass</strong> <strong>waveforms</strong> (using Prop. 3.2.1), all non-vanishing terms<br />

must be newforms. And by varying V , we see that the number <strong>of</strong> newforms is<br />

infinite.<br />

To use the above proposition we must find a permutation τ which conjugates<br />

the pair (σE,σR) into (σ ∗ E ,σ∗ R ). Then we must find an A ∈ PSL(2,Z) with<br />

h(A)(1) = τ(1). And finally we set α = JA and see if α2 ∈ G, and if so, we must<br />

also see how the various cusp classes transform under α. To be certain that<br />

we actually find newforms we must also investigate the possible supergroups<br />

between G and PSL(2,Z).<br />

There are two special cases, in which we can see the exact form <strong>of</strong> α and<br />

only have to investigate the permutation <strong>of</strong> the cusp classes and possible supergroups.<br />

Example 3.2.1. If we can take A as a power <strong>of</strong> S : A = S k . Then α : x + iy ↦→<br />

−x + k + iy, i.e. a reflection in the line x = k<br />

2 and it is clear that α2 = Id ∈ G.<br />

Example 3.2.2. If we can take A = E (i.e. we can do that if E(1) = 2 and<br />

σR ∼1 σ 2 R , which is always true if e2 = 0 and e3 ≥ 1). Then α = EJ and<br />

α : z ↦→ 1<br />

z , i.e. α is a reflection in the circle |z| = 1 and α2 = Id ∈ G.<br />

Remark 3.2.6. Consider a group G ⊆ PSL(2,Z), which has cusp width µ1 = L<br />

at ∞, and where J acts as a reflection in the line x = k<br />

2 , i.e. α : z ↦→ k − z.<br />

Suppose that f is a <strong>Maass</strong> waveform on G with α-eigenvalue ε = ±1, i.e.<br />

f (αz) = ε f (z). We know that the Fourier series <strong>of</strong> f at ∞ can be written as<br />

<br />

nx<br />

<br />

,<br />

L<br />

114<br />

f (z) = ∑ c(n)<br />

n=0<br />

√ yKiR(2π|n|y)e


where e(x) = e2πix . If we apply α we get<br />

f (αz) = ∑ c(n)<br />

n=0<br />

√ =<br />

<br />

2π|n|y n(−x + k)<br />

yKiR e<br />

L<br />

L<br />

<br />

nk √yKiR<br />

2π|n|y −nx<br />

∑ c(n)e<br />

e<br />

n=0 L<br />

L L<br />

<br />

−nk √yKiR<br />

2π|n|y<br />

<br />

nx<br />

<br />

e ,<br />

L<br />

L L<br />

= ∑ c(−n)e<br />

n=0<br />

hence if f (αz) = ε f (z) we must have c(−n) = εc(n)e <br />

nk<br />

L , for all n = 0, and<br />

in particular if we set c(1) = 1 we get<br />

<br />

e<br />

c(−1) =<br />

πi<br />

L (2k) , ε = 1, (the even case)<br />

e πi<br />

L (2k±L) , ε = −1, (the odd case).<br />

We can thus determine if a given cusp form f is <strong>of</strong> even or odd type with respect<br />

to α by looking at the c(−1) coefficient. Cf. Tables 3.5 and 3.6.<br />

We will now consider some explicit examples where we can take A = S k as<br />

in Example 3.2.1. See Table 3.1 for definition <strong>of</strong> these groups.<br />

Example 3.2.3. Consider the group G = Γ10,20. One can verify (by conjugating<br />

the corresponding permutations) that we can take A = Sk with k = 3, and<br />

α : z ↦→ −z + 3 is a reflection in the line x = 3<br />

2 . G has three cusps and we can<br />

choose representatives p1 = ∞, p2 = 1 and p3 = 3<br />

2 . Now we see that α(∞) = ∞,<br />

α(1) = 2 and α <br />

3 3<br />

2 = 2 . Hence, since α2 = Id also α(p2) = p2. And thus α<br />

preserves all cusp classes.<br />

For any α-odd newform f on G, the Fourier coefficient c(−1) = e πi<br />

4 (6±4) = i.<br />

Experimentally we found only odd newforms on this group! Cf. Remark 3.2.9<br />

and Table 3.6.<br />

Example 3.2.4. Consider the group G = Γ10,14. One can verify that here we<br />

can take A = S k with k = 6, and α : z ↦→ −z + 6 is a reflection in the line<br />

x = 3. G has three cusps and we can choose representatives p1 = ∞, p2 = 1<br />

and p3 = 3. Now we see that α(∞) = ∞, α(1) = 5 and α (3) = 3. Hence, since<br />

α 2 = Id also α(p2) ∼ p2. And thus α preserves all cusp classes.<br />

For any α-odd newform f on G, the Fourier coefficient c(−1) = e πi<br />

7 (12±7) =<br />

e 5πi<br />

7 . Experimentally we found only odd newforms on this group! Cf. Remark<br />

3.2.9 and Table 3.6.<br />

Example 3.2.5. Consider the group G = Γ10,8;7. One can verify that here we<br />

can take A = S k with k = 7. We get α : z ↦→ −z + 7. G has three cusps and<br />

we can choose representatives p1 ∈ {∞,0}, p2 = 2 and p3 = 5. We see that<br />

115


α(p1) = p1 and α (p2) = p3. Hence α permutes the cusps p2 and p3. And<br />

thus α does not preserve all cusp classes. In this case, experimentally, we<br />

found no <strong>Maass</strong> <strong>waveforms</strong> at all except for the oldforms from PSL(2,Z).<br />

Example 3.2.6. Consider the group G = Γ10,30. One can verify that here we<br />

can take A = S k with k = 4, and α : z ↦→ −z + 4 is a reflection in the line<br />

x = 2. G has three cusps and we can choose representatives p1 = ∞, p2 = 1<br />

and p3 = 2. We see that α(p1) = p1 and α(p3) = p3 and hence α (p2) = p2.<br />

Thus α preserves all cusp classes.<br />

For any α-odd newform f on G, the Fourier coefficient c(−1) = e πi<br />

5 (8±5) =<br />

e 3πi<br />

5 . Experimentally we found only odd newforms on this group! Cf. Remark<br />

3.2.9 and Table 3.6.<br />

Example 3.2.7. Consider the cycloidal group G = Γ10,10,1. One can verify that<br />

here we can take A = Sk with k = 9, and α : z ↦→ −z + 9 is a reflection in the<br />

line x = 9<br />

2 . G has only one cusp, ∞, and it is clearly preserved by α. Any<br />

α-odd function f on G will have Fourier coefficient c(−1) = e πi<br />

10 (18±10) = e 4πi<br />

5 ,<br />

and any α-even function will have c(−1) = e πi<br />

10 (18) = e 9πi<br />

5 . Experimentally we<br />

found both odd and even newforms on this group! Cf. Remark 3.2.9 and Table<br />

3.5.<br />

To ascertain that we actually have newforms we have to find a way to check<br />

for supergroups in PSL(2,Z); fortunately there is a well-known combinatorial<br />

method to find all these. We now turn to describe this method.<br />

Definition 3.2.3. A block system for a subgroup H ⊆ Sµ is an equivalence<br />

relation ≡ on the set {1,2,..., µ} such that<br />

∀h ∈ H, ∀x,y ∈ {1,2,..., µ} : x ≡ y ⇔ h(x) ≡ h(y).<br />

An equivalence class <strong>of</strong> ≡ is called a block.<br />

Note that for any H ⊆ Sµ there exist two trivial block systems: The block<br />

system with µ blocks, i.e. x ≡ y for all x = y ∈,{1,2,..., µ}; and the block<br />

system with 1 block, i.e. x ≡ y for all x,y ∈ {1,2,..., µ}.<br />

It is easy to see that if the group H is transitive, then all blocks <strong>of</strong> H must<br />

have the same number <strong>of</strong> elements. One can show the following proposition.<br />

Proposition 3.2.2. Let the legitimate pair (σE,σR) correspond to the group<br />

G ⊂ PSL(2,Z). Then there exists a supergroup G ′ , G ⊂ G ′ ⊂ PSL(2,Z) if and<br />

only if the group 〈σE,σR〉 = h(PSL(2,Z)) ⊆ Sµ supports a non-trivial block<br />

system.<br />

Outline <strong>of</strong> pro<strong>of</strong>. Suppose that G ′ is a supergroup <strong>of</strong> G in PSL(2,Z) and set<br />

B = {h(γ)(1)|γ ∈ G ′ } ⊆ {1,2,..., µ}. Observe that |B| ≥ 2 if G ′ = G, and<br />

116


|B| < µ if G ′ = PSL(2,Z). It is easy to verify that the equivalence relation<br />

defined by<br />

x ≡ y ⇐⇒ ∃γ ∈ PSL(2,Z) : [h(γ)(x) ∈ B and h(γ)(y) ∈ B]<br />

gives a block system, and that [1] = B.<br />

On the other hand, if ≡ is a non-trivial block system for h(PSL(2,Z)), then<br />

it is easy to verify that the group<br />

G ′ = {γ ∈ PSL(2,Z)|h(γ)(1) ≡ 1}<br />

is a supergroup <strong>of</strong> G, G G ′ PSL(2,Z).<br />

Remark 3.2.7. Note that an equivalence relation ≡ is a block system for the<br />

group h(PSL(2,Z)) = 〈σE,σR〉 if and only if<br />

∀x,y ∈ {1,..., µ} : x ≡ y ⇐⇒ σE(x) ≡ σE(y) ⇐⇒ σR(x) ≡ σR(y).<br />

Remark 3.2.8. Suppose that (abc) is a 3-cycle <strong>of</strong> σR. Then if a ≡ b, we must<br />

have σR(a) = b ≡ σR(b) = c, and similarly if a ≡ c or b ≡ c. Hence either<br />

a,b,c are in the same block or they are all in different blocks.<br />

We will now give an explicit example <strong>of</strong> how one checks for block systems.<br />

Example 3.2.8. Consider the group Γ10,20 (see Table 3.1). There are two possibilities<br />

for block systems. The block sizes are either 5+5 or 2+2+2+2+2.<br />

We have σE = (12)(34)(56)(78)(910) and σR = (1)(235)(467)(8910). By<br />

Remark 3.2.8 it is clear that we can not have a 5 + 5 block system. How about<br />

2 + 2 + 2 + 2 + 2? Suppose that 1 ≡ x, then 1 = σR(1) ≡ σR(x), hence the<br />

block containing 1 must contain x and σR(x). Since the block size is 2 we must<br />

have σR(x) = x or σR(x) = 1, but the only fixed element <strong>of</strong> σR is 1, hence it is<br />

impossible to have a 2 + 2 + 2 + 2 + 2 block system also.<br />

We conclude that there is no supergroup between Γ10,20 and PSL(2,Z). And<br />

in view <strong>of</strong> Proposition 3.2.1, Example 3.2.3 and Remark 3.2.5 we have shown<br />

that there exist newforms on Γ10,20.<br />

Remark 3.2.9. The groups we are primarily interested in are those which are<br />

not congruence subgroups and that satisfy Prop. 3.2.1. Up to index 11 (which<br />

is as far we have checked) this includes all groups G which satisfy (G ∗ ) = (G)<br />

except for Γ10,8,3 (on which the symmetry α permutes the cusps). All these<br />

groups have been investigated with respect to possible supergroups (note that<br />

there can be no supergroup <strong>of</strong> groups with prime index).<br />

Of the groups under consideration we found no supergroups for: all groups<br />

<strong>of</strong> index 10, the groups Γ9,9;1, Γ9,20, Γ9,15 and Γ9,15 <strong>of</strong> index 9, Γ8,8;1 and Γ8,8;3<br />

<strong>of</strong> index 8. On all these groups we claim that there exist newforms.<br />

117


3.2.2 Congruence subgroups<br />

Let G ⊆ PSL(2,Z) be a subgroup <strong>of</strong> index µ and generalized level N. It is<br />

known that if G is a congruence subgroup then G has level N (cf. [104]).<br />

We also remark that any PSL(2,Z)-conjugate <strong>of</strong> a congruence group is also a<br />

congruence subgroup since Γ(N) is normal in PSL(2,Z).<br />

Remark 3.2.10. Observe that if G has index µ and is a congruence subgroup<br />

<strong>of</strong> level N, then µ must divide the index <strong>of</strong> Γ(N) in Γ(1) = PSL(2,Z):<br />

[Γ(N) : Γ(1)] = 1<br />

2 N3 −2<br />

∏ 1 − p<br />

p|N<br />

, (3.3)<br />

where the product is taken over all prime divisors <strong>of</strong> N (valid for N ≥ 3, for<br />

N = 2 drop the factor 1<br />

2 ). Hence if this is not the case, then G can not be a<br />

congruence subgroup. The definitions <strong>of</strong> the groups that appear in the next two<br />

examples are given in Table 3.1.<br />

Example 3.2.9. Consider the groups Γ7,6, Γ7,10and Γ7,12 <strong>of</strong> type (7; 0,2,1,1)<br />

and generalized levels 6, 10 and 12 respectively. From (3.3) we see that the<br />

index <strong>of</strong> Γ(6) in PSL(2,Z) is 72, and the index <strong>of</strong> Γ(10) in PSL(2,Z) is 360.<br />

Hence since 7 ∤ 72 and 7 ∤ 360, all the above groups are noncongruence subgroups.<br />

Example 3.2.10. Consider the group Γ9,8 <strong>of</strong> type (9; 0,2,3,0) and generalized<br />

level 8. The index <strong>of</strong> Γ(8) in Γ(1) is 192, and since 9 ∤ 192 the group Γ9,8 is a<br />

noncongruence subgroup.<br />

Clearly this fact can not tell us that a given group is a congruence subgroup,<br />

but it is even the case that in general it can not even help us to determine<br />

that a given group is not a congruence subgroup. E.g. the group Γ9,6;1 <strong>of</strong><br />

type (9; 0,2,3,0) and generalized level 6 is not a congruence subgroup even<br />

though 9|72. But actually the group Γ9,6;2 with the same type is a congruence<br />

subgroup. So the type and generalized level are not enough to decide if a<br />

group is a congruence group or not.<br />

Fortunately there is a stronger theorem that give us necessary and sufficient<br />

conditions on the legitimate pair (σE,σR) in order for the corresponding subgroup<br />

to be a congruence subgroup. The following theorem is reformulated<br />

using our notation (and remember that we use an anti-homomorphism).<br />

Theorem 3.4. (Hsu [47, thm. 3.1])<br />

Suppose that (σE,σR) is a legitimate pair with corresponding group Γ and<br />

that σS = σRσE has order N = 2 k m, where m is odd and k ≥ 0. Let a and b<br />

be inverses <strong>of</strong> 2 and 5 mod N respectively (i.e. 2a ≡ 5b ≡ 1 mod N). We now<br />

have three cases:<br />

118


(a) N is odd, i.e. k = 0. Then Γ is a congruence subgroup if and only<br />

if<br />

1 = σRσSσRσ 1−a3<br />

S .<br />

(b) N is a power <strong>of</strong> 2, i.e. m = 1 : Let σS ˜ = σ 20<br />

S (σRσS) b σ −5<br />

S<br />

σ −1<br />

R .<br />

Then Γ is a congruence subgroup if and only if the following three<br />

relations hold:<br />

σ −1<br />

E ˜σSσE = ˜σ −1<br />

S , (b:1)<br />

˜σ −1<br />

S (σRσS) ˜σS = (σRσS) 25 , (b:2)<br />

<br />

˜σS (σRσS) 5 3 σE<br />

= 1. (b:3)<br />

(c) N has both an even and an odd part, i.e. if m ≥ 3 and k ≥ 1.<br />

Let c be the unique integer mod N such that c ≡ 0 mod 2 k and<br />

c ≡ 1 mod m, and let d be the unique integer mod N such that<br />

d ≡ 0 mod m and d ≡ 1 mod 2 k . Let α = σ c S , β = (σRσS) c , γ = σ d,<br />

S ,<br />

ρ = (σRσS) d, , π = γ 20 ρ b γ −4 ρ −1 , ζ = αβ −1 α and ξ = γρ −1 γ.<br />

Then Γ is a congruence subgroup if and only if the following 7<br />

relations hold:<br />

3.3 Numerical implementation<br />

αρ = ρα, (c:1)<br />

ζ 4 = 1, (c:2)<br />

ζ 2 = β −1 α 3 , (c:3)<br />

ζ 2 = β 2 α −a 3 , (c:4)<br />

ξ 2 = πρ 5 ξ 3 , (c:5)<br />

ξ −1 πξ = π −1 , (c:6)<br />

π −1 ρπ = ρ 25, . (c:7)<br />

Our goal is to investigate <strong>Maass</strong> <strong>waveforms</strong> on non-congruence subgroups,<br />

and thus we must find examples <strong>of</strong> such groups.<br />

Let N be a positive integer. By Theorem 3.2 we see that to find all subgroups<br />

(up to conjugation) <strong>of</strong> index N in PSL(2,Z) we first need to find all equivalence<br />

classes <strong>of</strong> legitimate pairs (σE,σR) in SN. To generate all elements <strong>of</strong><br />

SN is numerically speedy, and there exist a number <strong>of</strong> fast algorithms (cf.,<br />

e.g. [49]).<br />

The computationally hard part here is to identify the conjugates <strong>of</strong> legiti-<br />

119


mate pairs, i.e. given a list <strong>of</strong> legitimate pairs, find out the ones which are conjugates<br />

<strong>of</strong> each other. We know that conjugation preserves the cycle-structure,<br />

i.e. if σ1 = σ −1 σ2σ then σ1 and σ2 have the same set <strong>of</strong> cycles. If σ conjugates<br />

two pairs ( mod 1) we know first <strong>of</strong> all that σ(1) = 1, hence there are<br />

only (N − 1)! choices, and we can narrow the number <strong>of</strong> possibilities down<br />

further by studying the cycle structures <strong>of</strong> σE and σR. Unfortunately it seems<br />

that without calculating the details by hand for a specific value <strong>of</strong> N the problem<br />

is still hard even with as much automated simplifications as possible.<br />

Once we have a list <strong>of</strong> non-conjugated legitimate pairs we use Theorem<br />

3.4 to tell which pairs correspond to congruence subgroups. Then in order<br />

to apply the usual <strong>Maass</strong> waveform algorithms (cf. Section 1.3) for the corresponding<br />

subgroup <strong>of</strong> PSL(2,Z) we use Remark 3.2.1 to compute a set <strong>of</strong><br />

coset representatives.<br />

Example 3.3.1. We will see how to explicitly compute representatives for the<br />

PSL(2,Z)-conjugacy classes <strong>of</strong> the subgroups <strong>of</strong> type (7;0;1,1). First we observe<br />

that σE consists <strong>of</strong> one fixed element and 3 2-cycles. The letter 1 is<br />

special due to our convention in Theorem 3.2, but since we are only interested<br />

in representatives we are free to conjugate the fixed element <strong>of</strong> σE to 1 (cf. Remark<br />

3.2.2), hence, without loss <strong>of</strong> generality, we can assume that<br />

σE = (1)(23)(45)(67).<br />

Remember that the equivalence we consider (∼1) amounts to finding a permutation<br />

π which preserves 1 and now that we have fixed σE, such π must<br />

also preserve σE. Any choices we now make for σR are made modulo such a<br />

permutation.<br />

Now, we know that σR has one fixed element and 2 3-cycles. Since the<br />

numbers 2,3,...,7 are now on equal footing we can assume that σR(2) = 2<br />

and we can write<br />

σR = (2)(1bc)(de f ),<br />

and we have to choose the 3-cycles in turn. The elements 4,5,6,7 are equivalent,<br />

while 3 is essentially different. There are thus 2 cases for b. Either b = 3<br />

or b = 4. If b = 3, then we can assume c = 4 (and without loss <strong>of</strong> generality<br />

d = 5) and we have<br />

= (2)(134)(5e f ).<br />

σR1<br />

In case b = 4, then 6,7 are equivalent, and we observe that since (σE,σR) must<br />

be transitive we can not have c = 5. Hence either c = 3, or c = 6 (and, wlog<br />

d = 5 in both cases). We have<br />

120<br />

σR2<br />

σR3<br />

= (2)(143)(5e f ), and<br />

= (2)(146)(3e f ).


We now have to choose the last 3-cycle. For the last cycle <strong>of</strong> σR1 and σR2 ,<br />

clearly 6 and 7 are equivalent, and wlog e = 6 and f = 7 and we get<br />

σR1<br />

σR2<br />

= (2)(134)(567), and<br />

= (2)(143)(567).<br />

For σR3 = (2)(146)(5e f ), on the other hand, there are two possible (inequivalent)<br />

choices e = 5, f = 7 or e = 7, f = 5, corresponding to<br />

σR3<br />

σR4<br />

= (2)(146)(357), and<br />

= (2)(146)(375).<br />

We thus have 4 inequivalent legitimate pairs: σE,σR j<br />

subgroups G j, j = 1,2,3,4. The corresponding σS j = σR j σE are<br />

σS1<br />

σS2<br />

σS3<br />

σS4<br />

= (132465)(7),<br />

= (146532)(7),<br />

= (147)(2563),<br />

= (14327)(56).<br />

, corresponding to 4<br />

And we see that the generalized levels are 6, 6, 12 and 10. Note that σ ∗ R1 =<br />

σEσ 2 R1 σE = π −1 σR2 π ∼1 σR2 where π = (32)(45). Hence G∗ 1 = G2. The corresponding<br />

groups in Table 3.1 are Γ7,6, Γ ∗ 7,6 , Γ7,12 and Γ7,10. Fundamental<br />

domains <strong>of</strong> these groups, with the combinatorial structure indicated, can be<br />

seen in Figure 3.1.<br />

3.4 Numerical Results<br />

3.4.1 Subgroups<br />

In this section we will give examples <strong>of</strong> results obtained. Table 3.1 contains<br />

lists <strong>of</strong> specifications <strong>of</strong> conjugacy classes <strong>of</strong> subgroups <strong>of</strong> PSL(2,Z) <strong>of</strong> index<br />

less than or equal to 12. For each conjugacy class we have chosen a representative<br />

group, Γµ,l;i, for which we have also given the corresponding permutations.<br />

Here µ is the index, l is the generalized level and i is an arbitrary (but<br />

fixed by the program) enumeration <strong>of</strong> the conjugacy classes.<br />

To show which groups have a reflectional symmetry (cf. Prop. 3.2.1) we<br />

also list the equivalence class <strong>of</strong> G ∗ = JGJ. In this column we use the notation<br />

(G ∗ ) to say that (G ∗ ) = (G), and (G) to say that (G ∗ ) = (G) for the given<br />

group G.<br />

121


Figure 3.1: Fundamental domains for groups <strong>of</strong> type (7;0,1,1)<br />

122<br />

1 3 2 4 6 5<br />

(a) Γ7,6<br />

1 4 6 5 3 2<br />

1 4 7<br />

6<br />

3<br />

5<br />

2<br />

(c) Γ7,12<br />

7<br />

(b) Γ ∗ 7,6<br />

1 4 3 2 6<br />

7<br />

5<br />

(d) Γ7,10<br />

7


3.4.2 Newforms<br />

We have checked for <strong>Maass</strong> <strong>waveforms</strong> (newforms) on all the listed groups<br />

with index ≤ 11. As expected we found newforms on all congruence subgroups<br />

and all cycloidal subgroups. 1<br />

We also found newforms on all noncongruence, noncycloidal subgroups<br />

which satisfy Proposition 3.2.1. In this case, the number <strong>of</strong> newforms that we<br />

found is approximately one half <strong>of</strong> what is indicated by Weyl’s law for the<br />

surface G\H, as is to be expected for the α-odd space L 2 − (G\H). (In this<br />

vein, cf. [88, pp. 448 (footnote 1), 670 (middle)].) Since our investigation has<br />

been directed at investigating as many groups as possible, we have generally<br />

not taken R beyond 10 or so for each group. (That the Weyl ratios are only<br />

approximately 50% is thus not very disturbing; taking R larger undoubtedly<br />

produces values much closer to 50%.)<br />

In the case <strong>of</strong> noncongruence, noncycloidal groups which do not satisfy<br />

Proposition 3.2.1, we found no newforms whatsoever. An example <strong>of</strong> a group<br />

satisfying (G ∗ ) = (G), but that violates the cusp class preserving condition<br />

is Γ10,8;3. This indicates that the Phillips-Sarnak conjecture is true even for<br />

subgroups <strong>of</strong> the modular group.<br />

When testing for existence/nonexistence <strong>of</strong> newforms we ran the program<br />

up to height R = 10 for indexes ≤ 10 and up to R = 5 for index 11.<br />

In all cases <strong>of</strong> non-congruence, non-cycloidal groups that we have checked,<br />

all newforms found have been <strong>of</strong> α−odd type. For cycloidal groups, we found<br />

both α-even and α-odd newforms, exactly as one would expect from the trace<br />

formula. Cf. footnote 1.<br />

Tables 3.2-3.6 contain examples <strong>of</strong> eigenvalues for various noncongruence<br />

subgroups.<br />

Further data than that which is presented in this section is available from<br />

the author upon request.<br />

3.4.3 Fourier coefficients<br />

We have also investigated the distribution <strong>of</strong> Fourier coefficients in a few examples.<br />

Figure 3.3 (b) gives a histogram <strong>of</strong> 12422 Fourier coefficients (they<br />

are real and we normalized them to have mean 0 and standard deviation 1) <strong>of</strong> a<br />

newform with eigenvalue R = 3.34168098977 on Γ7,12. In the same figure we<br />

also plotted the Gaussian curve with the same mean and standard deviation.<br />

We have also studied the Fourier coefficients <strong>of</strong> a newform on Γ10,8;1 with<br />

eigenvalue R = 2.40931877484576. Here one should notice that the symmetry<br />

J acts as a shifted reflection on Γ10,8;1 (i.e. Γ ∗ 10,8;1 ∼ S−5 Γ10,8;1S 5 ) and hence<br />

the Fourier coefficients will not be real in this case, but have arguments shifted<br />

1 The Eisenstein series and scattering determinant for cycloidal groups are readily seen to be<br />

essentially the same as for PSL(2,Z). Cf. [105] and see also [100].<br />

123


Figure 3.2: Plot <strong>of</strong> Fourier coefficients in the complex plane.<br />

2<br />

1.5<br />

1<br />

0.5<br />

0<br />

-0.5<br />

-1<br />

-1.5<br />

”14412 Coefficients for R = 2.0493 . . . on Γ10,8;1”<br />

-2<br />

-3 -2 -1 0 1 2 3 4<br />

by multiples π<br />

8 . The surprising thing here is that it seems that the different<br />

residue classes modulo 8 also have different standard deviation, as can be<br />

seen in figure 3.2, which shows a plot in the complex plane <strong>of</strong> the first 14412<br />

Fourier coefficients. The coefficients were shifted back to the real line and<br />

then normalized with appropriate standard deviation (according to the residue<br />

class n mod 8). Figure 3.3(a) contains a histogram and a Gaussian fit to the<br />

resulting sequence.<br />

In both <strong>of</strong> these cases the fit to the Gaussian looks good, but more data<br />

has to be analyzed before one can draw any conclusions. This is certainly<br />

something that needs further investigation.<br />

124


Figure 3.3: Distribution <strong>of</strong> Fourier Coefficients<br />

0.4<br />

0.35<br />

0.3<br />

0.25<br />

0.2<br />

0.15<br />

0.1<br />

0.05<br />

14412 Coefficients<br />

Γ 10,8;1<br />

R=2.0493<br />

0<br />

−4 −3 −2 −1 0 1 2 3 4 5<br />

0.4<br />

0.35<br />

0.3<br />

0.25<br />

0.2<br />

0.15<br />

0.1<br />

0.05<br />

(a) At the group Γ10,8;1<br />

12422 Coefficients<br />

Γ 7,12<br />

R=3.341...<br />

0<br />

−4 −3 −2 −1 0 1 2 3 4 5<br />

(b) At the group Γ7,12<br />

125


126<br />

Table 3.1: Conjugacy classes <strong>of</strong> subgroups <strong>of</strong> PSL(2,Z) with index ≤ 12<br />

Γµ,l,i h e2 e3 g G ∗ C a σE σR σS<br />

Index=1 (1 equivalence class)<br />

(Γ2,2) 1 0 2 0 (G) T (1 2) ( 1)( 2) ( 1 2)<br />

Index=3 (2 equivalence classes)<br />

(Γ3,3) 1 3 0 0 (G) T (1)(2)(3) (1 2 3) (1 2 3)<br />

(Γ3,2) 2 1 0 0 (G) T (1)(2 3) (1 2 3) (1 2)(3)<br />

Index=4 (2 equivalence classes)<br />

(Γ4,4) 1 2 1 0 (G) T (1)(2)(3 4) (1 2 4)(3) (1 2 4 3)<br />

(Γ4,3) 2 0 1 0 (G) T (1 2)(3 4) (1)(2 3 4) (1 3 2)(4)<br />

Index=5 (1 equivalence class)<br />

(Γ5,5) 1 1 2 0 (G) T (1)(2 3)(4 5) (1 3 5)(2)(4) (1 3 2 5 4)<br />

Index=6 (10 equivalence classes)<br />

(Γ6,6;1) 1 0 0 1 (G) T (1 2)(3 4)(5 6) (1 3 5)(2 4 6) (1 4 5 2 3 6)<br />

(Γ6,6;2) 1 0 3 0 (G) T (1 2)(3 4)(5 6) (1)(2 4 6)(3)(5) (1 4 3 6 5 2)<br />

(Γ6,6;3) 1 4 0 0 (G) T (1)(2)(3)(4)(5 6) (1 2 5)(3 4 6) (1 2 5 3 4 6)<br />

(Γ6,5) 2 2 0 0 (G) T (1)(2)(3 4)(5 6) (1 2 3)(4 5 6) (1 2 3 5 4)(6)<br />

(Γ6,3) 2 2 0 0 (G) T (1)(2)(3 4)(5 6) (1 3 5)(2 4 6) (1 3 6)(2 4 5)<br />

a This column shows if the listed group is a congruence subgroup (T) or not (F).


127<br />

Γµ,l,i h e2 e3 g G ∗ C a σE σR σS<br />

(Γ6,4;1) 2 2 0 0 (G) T (1)(2)(3 4)(5 6) (1 3 6)(2 5 4) (1 3 2 5)(4 6)<br />

(Γ6,4;2) 3 0 0 0 (G) T (1 2)(3 4)(5 6) (1 2 3)(4 5 6) (1 3 5 4)(2)(6)<br />

(Γ6,4;3) 3 0 0 0 (G) T (1 2)(3 4)(5 6) (1 3 2)(4 6 5) (1)( 2 3 6 4)(5)<br />

(Γ6,4;4) 3 0 0 0 (G) T (1 2)(3 4)(5 6) (1 4 3)(2 6 5) (1 6 2 4)(3)(5)<br />

(Γ6,2) 3 0 0 0 (G) T (1 2)(3 4)(5 6) (1 3 6)(2 5 4) (1 5)(2 3)(4 6)<br />

Index=7 (6 equivalence classes)<br />

(Γ7,7) 1 3 1 0 (G ∗ ) T (1)(2)(3)(4 5)(6 7) (1 2 6)(3 5 7)(4) (1 2 6 3 5 4 7)<br />

(Γ7,6) 2 1 1 0 (G ∗ ) F (1)(2 3)(4 5)(6 7) (1 3 4)(2)(5 6 7) (1 3 2 4 6 5)(7)<br />

(Γ7,12) 2 1 1 0 (G) F (1)(2 3)(4 5)(6 7) (1 4 6)(2)(3 5 7) (1 4 7)(2 5 6 3)<br />

(Γ7,10) 2 1 1 0 (G) F (1)(2 3)(4 5)(6 7) (1 4 7)(2)(3 6 5) (1 4 3 2 6)(5 7)<br />

Index=8 (6 equivalence classes)<br />

(Γ8,8;1) 1 2 2 0 (G) F (1)(2)(3 4)(5 6)(7 8) (1 2 7)(3)(4 6 8)(5) (1 2 7 4 3 6 5 8)<br />

(Γ8,8;2) 1 2 2 0 (G ∗ ) T (1)(2)(3 4)(5 6)(7 8) (1 4 7)(2 6 8)(3)(5) (1 4 3 7 2 6 5 8)<br />

(Γ8,8;3) 1 2 2 0 (G) F (1)(2)(3 4)(5 6)(7 8) (1 4 8)(2 7 6)(3)(5) (1 4 3 8 6 5 2 7)<br />

(Γ8,7) 2 0 2 0 (G) T (1 2)(3 4)(5 6)(7 8) (1)(2 4 5)(3)(6 7 8) (1 4 3 5 7 6 2)(8)<br />

(Γ8,4) 2 0 2 0 (G) T (1 2)(3 4)(5 6)(7 8) (1)(2 5 7)(3)(4 6 8) (1 5 8 2)(3 6 7 4)<br />

(Γ8,6) 2 0 2 0 (G) T (1 2)(3 4)(5 6)(7 8) (1)(2 5 8)(3)(4 7 6) (1 5 4 3 7 2)(6 8)<br />

a This column shows if the listed group is a congruence subgroup (T) or not (F).


128<br />

Γµ,l,i h e2 e3 g G ∗ C a σE σR σS<br />

Index=9 (14 equivalence classes)<br />

(Γ9,9;1) 1 1 0 1 (G) F (1)(2 3)(4 5)(6 7)(8 9) (1 2 4)(3 6 8)(5 7 9) (1 2 6 9 3 4 7 8 5)<br />

(Γ9,9;2) 1 1 3 0 (G ∗ ) F (1)(2 3)(4 5)(6 7)(8 9) (1 3 8)(2)(4)(5 7 9)(6) (1 3 2 8 5 4 76 9)<br />

(Γ9,9;3) 1 5 0 0 (G) T (1)(2)(3)(4)(5)(6 7)(8 9) (1 2 6)(3 4 8)(5 7 9) (1 2 6 9 3 4 8 5 7)<br />

(Γ9,8) 2 3 0 0 (G ∗ ) F (1)(2)(3)(4 5)(6 7)(8 9) (1 2 4)(3 5 6)(7 8 9) (1 2 4 6 8 7 3 5)(9)<br />

(Γ9,6;1) 2 3 0 0 (G) F (1)(2)(3)(4 5)(6 7)(8 9) (1 2 4)(3 6 8)(5 7 9) (1 2 4 7 8 5)(36 9)<br />

(Γ9,14) 2 3 0 0 (G) F (1)(2)(3)(4 5)(6 7)(8 9) (1 2 4)(3 6 8)(5 9 7) (1 2 4 9 3 6 5)(7 8)<br />

(Γ9,20) 2 3 0 0 (G) F (1)(2)(3)(4 5)(6 7)(8 9) (1 4 6)(2 5 8)(3 7 9) (1 4 8 3 7)(2 56 9)<br />

(Γ9,6;2) 2 3 0 0 (G) T (1)(2)(3)(4 5)(6 7)(8 9) (1 4 6)(2 7 9)(3 8 5) (1 4 3 8 2 7)(5 6 9)<br />

(Γ9,15) 3 1 0 0 (G) F (1)(2 3)(4 5)(6 7)(8 9) (1 2 4)(3 5 6)(7 8 9) (1 2 5)(3 4 6 8 7)(9)<br />

(Γ9,6;3) 3 1 0 0 (G) F (1)(2 3)(4 5)(6 7)(8 9) (1 2 4)(3 6 5)(7 8 9) (1 2 6 8 7 5)(3 4)(9)<br />

(Γ9,7) 3 1 0 0 (G) F (1)(2 3)(4 5)(6 7)(8 9) (1 2 4)(3 6 7)(5 8 9) (1 2 6 3 4 8 5)(7)(9)<br />

(Γ9,12) 3 1 0 0 (G) F (1)(2 3)(4 5)(6 7)(8 9) (1 2 4)(3 6 8)(5 9 7) (1 2 6 5)(3 4 9)(7 8)<br />

Index=10 (26 equivalence classes)<br />

(Γ10,10;1) 1 0 1 1 (G) F (1 2)(3 4)(5 6)(7 8)(9 10) (1)(2 3 5)(4 7 9)(6 8 10) (1 3 7 10 4 5 8 9 6 2)<br />

(Γ10,10;2) 1 0 4 0 (G) T (1 2)(3 4)(5 6)(7 8)(9 10) (1)(2 4 9)(3)(5)(6 8 10)(7) (1 4 3 9 6 5 8 7 10 2)<br />

(Γ10,10;3) 1 4 1 0 (G) F (1)(2)(3)(4)(5 6)(7 8)(9 10) (1 2 7)(3 4 9)(5)(6 8 10) (1 2 7 10 3 4 9 6 5 8)<br />

a This column shows if the listed group is a congruence subgroup (T) or not (F).


129<br />

Γµ,l,i h e2 e3 g G ∗ C a σE σR σS<br />

(Γ10,10;4) 1 4 1 0 (G ∗ ) F (1)(2)(3)(4)(5 6)(7 8)(9 10) (1 2 7)(3 6 9)(4 8 10)(5) (1 2 7 10 3 6 5 9 4 8)<br />

(Γ10,10;5) 1 4 1 0 (G ∗ ) F (1)(2)(3)(4)(5 6)(7 8)(9 10) (1 2 7)(3 6 10)(4 9 8)(5) (1 2 7 4 9 3 6 5 10 8)<br />

(Γ10,9;1) 2 2 1 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10) (1 2 5)(3)(4 6 7)(8 9 10) (1 2 5 7 9 8 4 3 6)(10)<br />

(Γ10,9;2) 2 2 1 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10) (1 4 7)(2 5 8)(3)(6 9 10) (1 4 3 7 2 5 9 6 8)(10)<br />

(Γ10,9;3) 2 2 1 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10) (1 4 5)(2 6 7)(3)(8 9 10) (1 4 3 5 7 9 8 2 6)(10)<br />

(Γ10,12;1) 2 2 1 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10) (1 2 5)(3)(4 7 9)(6 8 10) (1 2 5 8 9 6)(3 7 10 4)<br />

(Γ10,12;2) 2 2 1 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10) (1 5 7)(2 6 9)(3)(48 10) (1 5 9 4 3 8)(2 6 7 10)<br />

(Γ10,8;1) 2 2 1 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10) (1 2 5)(3)(4 7 10)(6 9 8) (1 2 5 9 4 3 7 6)(8 10)<br />

(Γ10,8;2) 2 2 1 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10) (1 4 7)(2 5 9)(3)(6 8 10) (1 4 3 7 10 2 5 8)(6 9)<br />

(Γ10,21;1) 2 2 1 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10) (1 4 7)(2 5 9)(3)(6 10 8) (1 4 3 7 6 9 8)(2 5 10)<br />

(Γ10,21;2) 2 2 1 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10) (1 5 7)(2 8 10)(3)(49 6) (1 5 4 3 9 2 8)(6 7 10)<br />

(Γ10,5) 2 2 1 0 (G) T (1)(2)(3 4)(5 6)(7 8)(9 10) (1 5 7)(2 8 10)(3)(4 6 9) (1 5 9 2 8)(3 6 7 10 4)<br />

(Γ10,20) 3 0 1 0 (G) F (1 2)(3 4)(5 6)(7 8)(9 10) (1)(2 3 5)(4 6 7)(8 9 10) (1 3 6 2)(4 5 7 9 8)(10)<br />

(Γ10,14) 3 0 1 0 (G) F (1 2)(3 4)(5 6)(7 8)(9 10) (1)(2 3 5)(4 7 6)(8 9 10) (1 3 7 9 8 6 2)(4 5)(10)<br />

(Γ10,8;3) 3 0 1 0 (G) a F (1 2)(3 4)(5 6)(7 8)(9 10) (1)(2 3 5)(4 7 8)(6 9 10) (1 3 7 4 5 9 6 2)(8)(10)<br />

(Γ10,30) 3 0 1 0 (G) F (1 2)(3 4)(5 6)(7 8)(9 10) (1)(2 3 5)(4 7 10)(6 9 8) (1 3 7 6 2)(4 5 9)(8 10)<br />

Index=11 (26 equivalence classes)<br />

a This column shows if the listed group is a congruence subgroup (T) or not (F).


130<br />

Γµ,l,i h e2 e3 g G ∗ C a σE σR σS<br />

(Γ11,11;1) 1 3 2 0 (G) F (1)(2)(3)(4 5)(6 7)(8 9)(10 11) (1 8 10)(2 9 7)(3 5 11)(4)(6) (1 8 7 6 2 9 10 3 5 4 11)<br />

(Γ11,11;2) 1 3 2 0 (G) F (1)(2)(3)(4 5)(6 7)(8 9)(10 11) (1 5 10)(2 11 8)(3 9 7)(4)(6) (1 5 4 10 8 7 6 3 9 2 11)<br />

(Γ11,11;3) 1 3 2 0 (G ∗ ) F (1)(2)(3)(4 5)(6 7)(8 9)(10 11) (1 5 10)(2 9 3)(4)(6)(7 11 8) (1 5 4 10 8 3 2 9 7 6 11)<br />

(Γ11,11;4) 1 3 2 0 (G ∗ ) T (1)(2)(3)(4 5)(6 7)(8 9)(10 11) (1 2 10)(3 5 9)(4)(6)(7 11 8) (1 2 10 8 3 5 4 9 7 6 11)<br />

(Γ11,11;5) 1 3 2 0 (G ∗ ) F (1)(2)(3)(4 5)(6 7)(8 9)(10 11) (1 2 8)(3 9 11)(4)(5 10 7)(6) (1 2 8 11 7 6 5 4 10 3 9)<br />

(Γ11,11;6) 1 3 2 0 (G ∗ ) F (1)(2)(3)(4 5)(6 7)(8 9)(10 11) (1 8 10)(2 11 7)(3 9 5)(4)(6) (1 8 5 4 3 9 10 7 6 2 11)<br />

(Γ11,28;1) 2 1 2 0 (G) F (1)(2 3)(4 5)(6 7)(8 9)(10 11) (1 8 6)(2)(3 9 11)(4)(5 10 7) (1 8 11 7)(2 9 6 5 4 10 3)<br />

(Γ11,28;2) 2 1 2 0 (G ∗ ) F (1)(2 3)(4 5)(6 7)(8 9)(10 11) (1 6 3)(2)(4)(5 10 8)(7 11 9) (1 6 11 8 7 3 2)(4 10 9 5)<br />

(Γ11,10;1) 2 1 2 0 (G ∗ ) F (1)(2 3)(4 5)(6 7)(8 9)(10 11) (1 3 6)(2)(4)(5 7 11)(8 10 9) (1 3 2 6 11 9 10 5 4 7)(8)<br />

(Γ11,10;2) 2 1 2 0 (G ∗ ) F (1)(2 3)(4 5)(6 7)(8 9)(10 11) (1 6 3)(2)(4)(5 7 11)(8 10 9) (1 6 11 9 10 5 4 7 3 2)(8)<br />

(Γ11,10;3) 2 1 2 0 (G ∗ ) F (1)(2 3)(4 5)(6 7)(8 9)(10 11) (1 6 8)(2)(3 7 5)(4)(9 11 10) (1 6 5 4 3 2 7 8 11 9)(10)<br />

(Γ11,18;1) 2 1 2 0 (G) F (1)(2 3)(4 5)(6 7)(8 9)(10 11) (1 8 6)(2)(3 5 11)(4)(7 9 10) (1 8 10 3 2 5 4 11 7)(6 9)<br />

(Γ11,18;2) 2 1 2 0 (G ∗ ) F (1)(2 3)(4 5)(6 7)(8 9)(10 11) (1 3 6)(2)(4)(5 9 11)(7 10 8) (1 3 2 6 10 5 4 9 7)(8 11)<br />

(Γ11,30) 2 1 2 0 (G ∗ ) F (1)(2 3)(4 5)(6 7)(8 9)(10 11) (1 6 8)(2)(3 9 10)(4)(5 7 11) (1 6 11 3 2 9)(4 7 8 10 5)<br />

(Γ11,24;1) 2 1 2 0 (G) F (1)(2 3)(4 5)(6 7)(8 9)(10 11) (1 8 6)(2)(3 5 10)(4)(7 11 9) (1 8 7)(2 5 4 10 9 6 11 3)<br />

(Γ11,24;2) 2 1 2 0 (G) F (1)(2 3)(4 5)(6 7)(8 9)(10 11) (1 6 8)(2)(3 9 11)(4)(5 10 7) (1 6 5 4 10 3 2 9)(7 8 11)<br />

Index=12 (81 equivalence classes)<br />

a This column shows if the listed group is a congruence subgroup (T) or not (F).


131<br />

Γµ,l,i h e2 e3 g G ∗ C a σE σR σS<br />

(Γ12,12;1) 1 2 0 1 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 2 3)(4 7 10)(5 8 11)(6 9 12) (1 2 3 7 11 6 8 10 12 5 9 4)<br />

(Γ12,12;2) 1 2 0 1 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 10 7)(4 11 8)(6 12 9) (1 3 11 9 7 4 5 12 8 2 10 6)<br />

(Γ12,12;3) 1 2 0 1 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 10 7)(4 12 9)(6 11 8) (1 3 12 8 2 10 4 5 11 9 7 6)<br />

(Γ12,12;4) 1 2 0 1 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 9 6)(4 12 8)(7 10 11) (1 3 12 7 4 5 2 9 11 8 10 6)<br />

(Γ12,12;6) 1 2 0 1 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 9 4)(6 12 8)(7 10 11) (1 3 2 9 11 8 10 4 5 12 7 6)<br />

(Γ12,12;7) 1 2 3 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 9 11)(2 10 8)(3)(4 6 12)(5)(7) (1 9 8 7 2 10 11 4 3 6 5 12)<br />

(Γ12,12;9) 1 2 3 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 9 11)(2 12 8)(3)(4 6 10)(5)(7) (1 9 4 3 6 5 10 11 8 7 2 12)<br />

(Γ12,12;11) 1 2 3 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 2 11)(3)(4 6 10)(5)(7)(8 12 9) (1 2 11 9 4 3 6 5 10 8 7 12)<br />

(Γ12,12;13) 1 2 3 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 4 11)(2 6 10)(3)(5)(7)(8 12 9) (1 4 3 11 9 2 6 5 10 8 7 12)<br />

(Γ12,12;15) 1 2 3 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 4 11)(2 10 6)(3)(5)(7)(8 12 9) (1 4 3 11 9 6 5 2 10 8 7 12)<br />

(Γ12,12;16) 1 2 3 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 4 11)(2 9 8)(3)(5)(6 10 12)(7) (1 4 3 11 6 5 10 8 7 2 9 12)<br />

(Γ12,12;17) 1 6 0 0 (G) F (1)(2)(3)(4)(5)(6)(7 8)(9 10)(11 12) (1 2 7)(3 12 6)(4 8 10)(5 9 11) (1 2 7 10 11 6 3 12 5 9 4 8)<br />

(Γ12,12;18) 1 6 0 0 (G ∗ ) F (1)(2)(3)(4)(5)(6)(7 8)(9 10)(11 12) (1 2 7)(3 811)(4 9 12)(5 10 6) (1 2 7 11 4 9 6 5 10 12 3 8)<br />

(Γ12,12;20) 1 6 0 0 (G) T (1)(2)(3)(4)(5)(6)(7 8)(9 10)(11 12) (1 2 7)(3 104)(5 12 6)(8 11 9) (1 2 7 11 6 5 12 9 4 3 10 8)<br />

(Γ12,8;1) 2 0 0 1 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 10 7)(4 12 9)(6 11 8) (1 10 4 5 11 9 7 6)(2 3 12 8)<br />

(Γ12,6;1) 2 0 0 1 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 11 9)(4 6 7)(8 12 10) (1 11 10 2 3 6)(4 5 7 12 9 8)<br />

(Γ12,10;1) 2 0 0 1 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 10 12)(4 6 7)(8 11 9) (1 10 8 4 5 7 11 2 3 6)(9 12)<br />

a This column shows if the listed group is a congruence subgroup (T) or not (F).


132<br />

Γµ,l,i h e2 e3 g G ∗ C a σE σR σS<br />

(Γ12,11;1) 2 0 0 1 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 9 6)(4 10 7)(8 11 12) (1 9 7 11 8 4 5 2 3 10 6)(12)<br />

(Γ12,9;1) 2 0 0 1 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 10 8)(4 11 9)(6 7 12) (1 10 4 5 7 2 3 11 6)(8 12 9)<br />

(Γ12,35;1) 2 0 3 0 (G) F (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1)(2 10 7)(3)(4 11 8)(5)(6 12 9) (1 10 6 5 12 8 2)(3 11 9 7 4)<br />

(Γ12,11;2) 2 0 3 0 (G) F (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1)(2 9 6)(3)(4 10 7)(5)(8 11 12) (1 9 7 11 8 4 3 10 6 5 2)(12)<br />

(Γ12,11;3) 2 0 3 0 (G) F (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1)(2 6 9)(3)(4 7 10)(5)(8 12 11) (1 6 5 9 4 3 7 12 8 10 2)(11)<br />

(Γ12,10;2) 2 0 3 0 (G) F (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1)(2 10 4)(3)(5)(6 11 8)(7 12 9) (1 10 7 6 5 11 9 4 3 2)(8 12)<br />

(Γ12,8;2) 2 0 3 0 (G) F (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1)(2 9 6)(3)(4 12 8)(5)(7 10 11) (1 9 11 8 10 6 5 2)(3 12 7 4)<br />

(Γ12,9;2) 2 0 3 0 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1)(2 8 10)(3)(4 9 11)(5)(6 12 7) (1 8 6 5 12 4 3 9 2)(7 10 11)<br />

(Γ12,11;4) 2 4 0 0 (G ∗ ) F (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 2 5)(3 6 8)(4 7 12)(9 11 10) (1 2 5 8 12 10 11 4 7 3 6)(9)<br />

(Γ12,9;3) 2 4 0 0 (G ∗ ) F (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 5 7)(2 11 9)(3 4 6)(8 12 10) (1 5 3 4 6 7 12 9 8)(2 11 10)<br />

(Γ12,8;3) 2 4 0 0 (G ∗ ) F (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 2 5)(3 10 7)(4 11 8)(6 12 9) (1 2 5 12 8 3 10 6)(4 11 9 7)<br />

(Γ12,35;2) 2 4 0 0 (G) F (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 2 5)(3 12 9)(4 10 7)(6 11 8) (1 2 5 11 9 7 6)(3 12 8 4 10)<br />

(Γ12,10;3) 2 4 0 0 (G ∗ ) F (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 5 7)(2 10 12)(3 4 6)(8 11 9) (1 5 3 4 6 7 11 2 10 8)(9 12)<br />

(Γ12,11;6) 2 4 0 0 (G ∗ ) F (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 5 7)(2 10 8)(3 4 6)(9 11 12) (1 5 3 4 6 7 2 10 11 9 8)(12)<br />

(Γ12,10;5) 2 4 0 0 (G) T (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 2 5)(3 10 4)(6 11 8)(7 12 9) (1 2 5 11 9 4 3 10 7 6)(8 12)<br />

(Γ12,35;3) 2 4 0 0 (G) F (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 5 7)(2 10 8)(3 6 12)(4 11 9) (1 5 12 9 8)(2 10 4 11 3 6 7)<br />

(Γ12,6;2) 2 4 0 0 (G) F (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 2 5)(3 11 4)(6 10 8)(7 12 9) (1 2 5 10 7 6)(3 11 9 8 12 4)<br />

a This column shows if the listed group is a congruence subgroup (T) or not (F).


133<br />

Γµ,l,i h e2 e3 g G ∗ C a σE σR σS<br />

(Γ12,6;3) 2 4 0 0 (G) T (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 5 7)(2 11 8)(3 12 9)(4 6 10) (1 5 10 3 12 8)(2 11 9 4 6 7)<br />

(Γ12,11;8) 2 4 0 0 (G) F (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 2 5)(3 4 7)(6 8 12)(9 11 10) (1 2 5 8 3 4 7 12 10 11 6)(9)<br />

(Γ12,9;5) 2 4 0 0 (G) F (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 2 5)(3 10 8)(4 11 9)(6 7 12) (1 2 5 7 3 10 4 11 6)(8 12 9)<br />

(Γ12,6;4) 2 4 0 0 (G) F (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 5 7)(2 12 6)(3 10 8)(4 11 9) (1 5 2 12 9 8)(3 10 4 11 6 7)<br />

(Γ12,8;5) 2 4 0 0 (G) T (1)(2)(3)(4)(5 6)(7 8)(9 10)(11 12) (1 5 7)(2 12 6)(3 8 10)(4 9 11) (1 5 2 12 4 9 3 8)(6 7 10 11)<br />

(Γ12,18;1) 3 2 0 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 2 3)(4 9 6)(5 10 7)(8 11 12) (1 2 3 9 7 11 8 5 4)(6 10)(12)<br />

(Γ12,30;1) 3 2 0 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 2 3)(4 7 9)(5 8 10)(6 12 11) (1 2 3 7 10 4)(5 12 6 8 9)(11)<br />

(Γ12,42;1) 3 2 0 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 2 3)(4 7 11)(5 8 10)(6 9 12) (1 2 3 7 10 12 4)(5 9)(6 8 11)<br />

(Γ12,6;5) 3 2 0 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 11 9)(4 6 7)(8 12 10) (1 3 6)(2 11 10)(4 5 7 12 9 8)<br />

(Γ12,12;21) 3 2 0 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 10 4)(6 11 8)(7 12 9) (1 3 2 10 7 6)(4 5 11 9)(8 12)<br />

(Γ12,42;2) 3 2 0 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 10 12)(4 6 7)(8 11 9) (1 3 6)(2 10 8 4 5 7 11)(9 12)<br />

(Γ12,24;1) 3 2 0 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 10 8)(4 6 7)(9 11 12) (1 3 6)(2 10 11 9 8 4 5 7)(12)<br />

(Γ12,28;1) 3 2 0 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 9 6)(4 10 7)(8 11 12) (1 3 10 6)(2 9 7 11 8 4 5)(12)<br />

(Γ12,4;1) 3 2 0 0 (G) T (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 11 8)(4 10 7)(6 12 9) (1 3 10 6)(2 11 9 7)(4 5 12 8)<br />

(Γ12,24;3) 3 2 0 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 9 4)(6 10 7)(8 11 12) (1 3 2 9 7 11 8 6)(4 5 10)(12)<br />

(Γ12,10;6) 3 2 0 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 2 3)(4 8 10)(5 9 6)(7 12 11) (1 2 3 8 12 7 10 6 9 4)(5)(11)<br />

(Γ12,60;1) 3 2 0 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 11 8)(4 12 9)(6 10 7) (1 3 12 8 6)(2 11 9 7)(4 5 10)<br />

a This column shows if the listed group is a congruence subgroup (T) or not (F).


134<br />

Γµ,l,i h e2 e3 g G ∗ C a σE σR σS<br />

(Γ12,10;7) 3 2 0 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 4 10)(6 11 8)(7 12 9) (1 3 10 7 6)(2 4 5 11 9)(8 12)<br />

(Γ12,30;2) 3 2 0 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 6 8)(4 7 12)(9 11 10) (1 3 7 2 6)(4 5 8 12 10 11)(9)<br />

(Γ12,10;9) 3 2 0 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 6 11)(4 7 8)(9 12 10) (1 3 7 4 5 11 10 12 2 6)(8)(9)<br />

(Γ12,10;10) 3 2 0 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 12 6)(4 7 8)(9 11 10) (1 3 7 4 5 2 12 10 11 6)(8)(9)<br />

(Γ12,6;6) 3 2 0 0 (G) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 11 9)(4 10 8)(6 7 12) (1 3 10 2 11 6)(4 5 7)(8 12 9)<br />

(Γ12,18;2) 3 2 0 0 (G ∗ ) F (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 11 9)(4 12 6)(7 10 8) (1 3 12 9 8 10 2 11 6)(4 5)(7)<br />

(Γ12,8;6) 3 2 0 0 (G) T (1)(2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 11 10)(4 8 6)(7 9 12) (1 3 8 9 2 11 7 6)(4 5)(10 12)<br />

(Γ12,6;7) 4 0 0 0 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 4)(2 7 12)(5 8 10)(6 9 11) (1 7 10 11 2 3)(4)(5 9)(6 8 12)<br />

(Γ12,8;7) 4 0 0 0 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 4)(2 9 6)(5 10 7)(8 11 12) (1 9 7 11 8 5 2 3)(4)(6 10)(12)<br />

(Γ12,9;6) 4 0 0 0 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 4)(2 8 5)(6 9 10)(7 12 11) (1 8 12 7 5 9 6 2 3)(4)(10)(11)<br />

(Γ12,5;1) 4 0 0 0 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 12 6)(4 7 8)(9 11 10) (1 12 10 11 6)(2 3 7 4 5)(8)(9)<br />

(Γ12,3;1) 4 0 0 0 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 11 9)(4 10 8)(67 12) (1 11 6)(2 3 10)(4 5 7)(8 12 9)<br />

(Γ12,6;8) 4 0 0 0 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 5 3)(2 9 11)(4 6 12)(7 8 10) (1 9 7 10 11 4)(2 5 12)(3 6)(8)<br />

(Γ12,4;2) 4 0 0 0 (G) T (1 2)(3 4)(5 6)(7 8)(9 10)(11 12) (1 3 5)(2 11 10)(4 8 6)(79 12) (1 11 7 6)(2 3 8 9)(4 5)(10 12)<br />

a This column shows if the listed group is a congruence subgroup (T) or not (F).


Table 3.2: Eigenvalues on Γ7,10 (2 cusps)<br />

R |c(−1) + 1| a h1<br />

2.98380890177954200 7E-12 2E-15<br />

4.25864509747714400 1E-12 4E-13<br />

4.89540660025545100 2E-13 4E-13<br />

5.28884863158980700 1E-12 6E-14<br />

6.15234838528998200 0E+00 2E-14<br />

6.27199072998424200 2E-13 4E-13<br />

6.76036517325738900 2E-13 2E-13<br />

7.20603858428835700 4E-13 1E-14<br />

7.55764967370570000 4E-13 7E-13<br />

7.77456459246736200 6E-14 3E-14<br />

8.50359685011489000 8E-13 8E-14<br />

8.77366479019266900 3E-13 2E-12<br />

8.88752983993304300 2E-13 2E-14<br />

9.29204073179663000 5E-14 6E-13<br />

9.53369526135354484 * 2E-14 3E-14<br />

9.67736151045037300 4E-13 6E-14<br />

10.04237676924058000 1E-10 2E-08<br />

All newforms listed are odd with respect to the symmetry α : z ↦→ −¯z + 5.<br />

* An oldform from PSL(2,Z).<br />

a See Remark 3.2.6 for an explanation <strong>of</strong> the term 1.<br />

135


Table 3.3: Eigenvalues on Γ9,20;1 (2 cusps)<br />

R |c(−1) − e 2πi<br />

10 | a h1<br />

2.05495976063979 6E-11 1E-15<br />

3.25563226951552 1E-12 2E-14<br />

4.01743421448455 8E-13 1E-12<br />

4.39990456866076 3E-13 9E-12<br />

4.92041181302557 1E-13 5E-14<br />

5.22354223285616 1E-11 6E-11<br />

5.70296517107618 2E-14 3E-13<br />

6.20783400944797 3E-12 6E-12<br />

6.52618304206566 8E-13 2E-12<br />

6.78783029888573 4E-13 5E-14<br />

7.07791461588707 2E-12 1E-14<br />

7.38877048107201 5E-15 8E-13<br />

8.02743455884534 3E-10 3E-09<br />

8.70656893384487 2E-13 7E-10<br />

8.81164911527479 6E-15 9E-13<br />

9.03359858422541 9E-11 3E-10<br />

9.50035839336283 2E-11 4E-11<br />

9.53369526135354 * 5E-09 7E-14<br />

9.66113248208174 3E-11 3E-12<br />

All newforms listed are odd with respect to the symmetry α : z ↦→ −¯z + 4.<br />

* An oldform from PSL(2,Z).<br />

a See Remark 3.2.6 for an explanation <strong>of</strong> the term e 2πi<br />

10 .<br />

136


Table 3.4: Eigenvalues <strong>of</strong> newforms on Γ9,15 (3 cusps)<br />

R |c(−1) + 1| a h1<br />

3.11274526127958 8E-13 2E-15<br />

3.79758314524666 8E-14 1E-14<br />

4.68277129730610 2E-14 3E-11<br />

5.61420134556306 9E-12 3E-11<br />

6.17840605444398 4E-14 1E-10<br />

6.28023134316972 2E-11 4E-10<br />

6.50085225008714 8E-10 2E-10<br />

6.87257413990397 8E-09 6E-10<br />

7.21783938535166 2E-10 9E-10<br />

7.41980107735753 2E-11 4E-11<br />

7.73811611339896 2E-12 4E-12<br />

7.89543101394190 3E-10 1E-10<br />

8.33429557635011 5E-11 6E-09<br />

8.82147607539715 4E-08 2E-10<br />

8.83035851126317 6E-09 2E-08<br />

9.17711548223084 3E-10 7E-10<br />

9.23041542661088 6E-08 3E-09<br />

9.53369526135299 * 4E-09 5E-11<br />

9.73675250028969 3E-09 3E-10<br />

All newforms listed are odd with respect to the symmetry α : z ↦→ −¯z + 3.<br />

* An oldform from PSL(2,Z).<br />

a See Remark 3.2.6 for an explanation <strong>of</strong> the term 1.<br />

137


Table 3.5: Eigenvalues <strong>of</strong> newforms on a cycloidal group <strong>of</strong> index 10<br />

Γ10,10,1<br />

R |c(−1) − e 9πi<br />

5 | a |c(−1) − e 4πi<br />

5 | a h1<br />

1.56460530796110 e 6E-11 3E-15<br />

2.11939343023923 o 1E-11 1E-12<br />

2.50414532509967 e 1E-09 2E-09<br />

2.91592762476298 e 6E-10 3E-14<br />

3.00023468927910 o 2E-11 4E-15<br />

3.10845111145684 o 1E-13 1E-11<br />

3.56237482722476 e 6E-11 2E-15<br />

3.65766022871834 o 1E-12 2E-13<br />

3.88274850431287 o 2E-12 5E-11<br />

4.10694094068138 e 6E-11 2E-14<br />

4.55947270613401 o 4E-12 4E-14<br />

4.74782424736486 o 3E-13 4E-15<br />

4.76752492536872 e 6E-11 2E-14<br />

4.84953817531720 e 1E-10 2E-12<br />

e Even with respect to the symmetry α : z ↦→ −¯z + 9.<br />

o Odd with respect to the symmetry α : z ↦→ −¯z + 9.<br />

a See Remark 3.2.6 for an explanation <strong>of</strong> the terms e 9πi<br />

5 and e 4πi<br />

5 .<br />

138


Table 3.6: Eigenvalues <strong>of</strong> newforms on groups <strong>of</strong> index 10<br />

Γ10,12,1 R |c(−1) − e 2πi<br />

3 | a h1<br />

2.78187318285282 7E-12 5E-16<br />

3.37184320841770 5E-11 7E-12<br />

3.76035106954282 3E-11 8E-12<br />

4.22985717336497 5E-13 6E-15<br />

4.85352368236249 3E-13 1E-12<br />

Γ10,8,1 R |c(−1) − e 3πi<br />

4 | a h1<br />

2.40931877484576 2E-10 3E-14<br />

3.29803900236809 5E-09 2E-10<br />

3.70953186018084 1E-11 3E-13<br />

4.27459659497848 2E-10 6E-11<br />

4.58646680301038 1E-12 2E-11<br />

Γ10,21,2 R |c(−1) − e 3πi<br />

7 | a h1<br />

1.90995590589148 5E-09 2E-14<br />

3.20456049497589 1E-08 6E-13<br />

3.50475869850690 5E-09 4E-12<br />

4.20753581320237 5E-09 6E-11<br />

4.77744371781533 5E-08 8E-10<br />

5.00937515698108 5E-09 3E-11<br />

Γ10,20 R |c(−1) − i| a h1<br />

3.08773394089329578 2E-10 1E-13<br />

3.50709021769279694 7E-10 2E-14<br />

4.34565944211539890 7E-10 3E-13<br />

Γ10,14 R |c(−1) − e 5πi<br />

7 | a h1<br />

2.57440388633114603 0E+00 2E-15<br />

3.79413397458154922 4E-11 2E-14<br />

4.06173226279882993 4E-11 2E-11<br />

4.68613704425140210 3E-11 1E-14<br />

Γ10,30 R |c(−1) − e 3πi<br />

5 | a h1<br />

2.68171983329039 2E-10 4E-16<br />

3.56451256590188 1E-10 4E-15<br />

4.44205124980089 7E-10 1E-12<br />

4.70189370885265 2E-10 1E-14<br />

4.97437106689348 2E-10 2E-15<br />

All newforms listed above are odd with respect to the action <strong>of</strong> J on the groups.<br />

a See Remark 3.2.6 for an explanation <strong>of</strong> the different terms e ·πi<br />

· .<br />

139


4 An Algorithm for Whittaker’s W-function<br />

4.1 Introduction<br />

We present a new algorithm for computing the Whittaker function Wl,iR(x)<br />

based on an integral representation similar to the integral representation <strong>of</strong> the<br />

K-Bessel function,<br />

KiR(x) = 1<br />

∞<br />

e<br />

2 −∞<br />

−xcosht · 1 · e iRt dt.<br />

In our case, the factor “1” in front <strong>of</strong> the exponential, e iRt , in the integrand is<br />

replaced with a confluent hypergeometric function, Ψ, with real parameters<br />

and complex argument. The advantage with this approach is that we avoid the<br />

computationally harder problem <strong>of</strong> evaluating the confluent hypergeometric<br />

functions with complex parameters. We then turn the path <strong>of</strong> the integral to<br />

a path <strong>of</strong> “almost stationary” phase, and the resulting integral turns out to be<br />

well-suited for standard quadrature techniques.<br />

It is clear that for this algorithm to be <strong>of</strong> any practical use we need efficient<br />

algorithms for the confluent hypergeometric function with real parameters.<br />

Fortunately there is an abundance <strong>of</strong> literature on this subject, and we will<br />

present three different methods.<br />

4.2 Presentation <strong>of</strong> the Algorithm<br />

We use the integral formula (cf. [39, p. 374]) to connect the confluent hypergeometric<br />

Ψ−function with complex parameters and real argument (occurring<br />

in the definition <strong>of</strong> the Whittaker W-function) to the same function with com-<br />

plex argument but real parameters:<br />

(2x) s− 1 2 e −x Ψ(s−l,2s;2x) = 1<br />

√ π<br />

∞<br />

0<br />

e −xcosht <br />

Ψ −l; 1<br />

2 ;ξ<br />

<br />

cosh s − 1<br />

<br />

t dt,<br />

2<br />

where ξ = x(1+cosht). The Whittaker W-function, Wl,iR(x), is the decreasing<br />

(as x → ∞) solution to Whittaker’s differential equation<br />

w ′′ <br />

+ − 1 l<br />

+<br />

4 x +<br />

1<br />

4 + R2<br />

x2 <br />

w = 0.<br />

141


We can express Wl,iR in terms <strong>of</strong> the Tricomi confluent hypergeometric function,<br />

Ψ, as (cf. [29, vol I, p. 264 (2)])<br />

Wl,iR(2x) = e −x (2x) 1 2 +iR <br />

1<br />

Ψ + iR − l, 1 + 2iR;2x .<br />

2<br />

With s = 1<br />

2<br />

+ iR we thus get the following integral representation:<br />

<br />

Wl,iR(2x) =<br />

−l; 1<br />

<br />

;x(1 + cosht) cosh[iRt]dt.<br />

2<br />

<br />

2x ∞<br />

e<br />

π 0<br />

−xcosht Ψ<br />

Hence the integral we want to evaluate is:<br />

∞<br />

I = e<br />

0<br />

−xcosht <br />

Ψ −l; 1<br />

<br />

;x(1 + cosht) cos(Rt)dt<br />

2<br />

= 1<br />

∞<br />

e<br />

2<br />

−xcosht <br />

Ψ −l; 1<br />

<br />

;x(1 + cosht) e<br />

2 iRt dt,<br />

that is we define<br />

−∞<br />

I = 1<br />

∞<br />

e<br />

2 −∞<br />

φ(t) ψ(t)dt, (4.1)<br />

where φ(t) = −xcosht + iRt and ψ(t) = Ψ −l; 1<br />

2 ;x(1 + cosht) .<br />

The trick is now to turn the contour out in the complex plane into a path <strong>of</strong><br />

“almost stationary” phase. The argument <strong>of</strong> ψ(t) is slowly varying compared<br />

to the argument <strong>of</strong> e φ(t) , hence a path <strong>of</strong> stationary phase for e φ(t) should<br />

provide a good approximation to the real path <strong>of</strong> stationary phase.<br />

The same integral, but with ψ(t) = e At has been studied extensively by<br />

Helen Avelin [10] in connection with the K-Bessel function and her notes<br />

have been very helpful.<br />

We now have to find the path <strong>of</strong> stationary phase for e φ(t) ; it will depend on<br />

the parameters x and R, or actually (as we will see) on the quotient<br />

4.3 The path<br />

T = R<br />

x .<br />

Along the desired path, the imaginary part <strong>of</strong> φ should be constant, and it<br />

should also pass through the point where the real part <strong>of</strong> φ has a maximum.<br />

This maximum is attained at the solutions <strong>of</strong><br />

142<br />

φ ′ (t) = −xsinht + iR = 0, (4.2)


which are the following points (depending on T ):<br />

<br />

t1 = ln(T +<br />

T ≥ 1 :<br />

√ T 2 − 1) + i π<br />

2 ,<br />

t2 = −ln(T + √ T 2 − 1) + i π<br />

2 ,<br />

T ≤ 1 : t = iarcsinT.<br />

We have to separate these two cases (T ≥ 1 and T ≤ 1), and for numerical<br />

reasons we will also treat the case <strong>of</strong> T ≈ 1 separately. We now view t as a<br />

complex variable, t = u + iv, and thus<br />

φ(t) = φ(u + iv)<br />

= −xcosh(u + iv) + iR(u + iv)<br />

= −x(coshucosv + isinhusinv) − Rv + iRu<br />

= −[xcoshucosv + Rv] − i[xsinhusinv − Ru].<br />

The desired path is determined by the equation<br />

ℑφ = Ru − sinhusinv = [constant], (4.3)<br />

where the constant must be the value <strong>of</strong> ℑφ at the solutions <strong>of</strong> (4.2). The<br />

path can be parametrized by u, hence we will sometimes use the notation<br />

φ(u) = φ(u + iv(u)) and ψ(u) = ψ(u + iv(u)). In the following sections we<br />

will detail the paths for the different cases <strong>of</strong> T .<br />

T ≥ 1 :<br />

There are two solutions <strong>of</strong> (4.2): t1 = u1 + iv1 and t2 = −u1 + iv1, where u1 =<br />

ln(T + √ T 2 − 1) and v1 = π<br />

2 . Hence we must split the path into two parts, but<br />

for numerical reasons we will actually split it into three parts.<br />

Consider first the critical point t1 = u1 + iv1. From (4.3) we see that<br />

ℑφ(u + iv) = ℑφ(v1)<br />

⇔<br />

Ru − xsinhusinv = Ru1 − xsinhu1 sinv1<br />

= Ru1 − xsinhu1 <br />

= Ru1 − x T 2 − 1,<br />

which, with S = Tu1 − √ T 2 − 1 = T ln(T + √ T 2 − 1)− √ T 2 − 1 is equivalent<br />

to<br />

Tu − S<br />

sinv =<br />

sinhu .<br />

143


This equation defines the path that goes through t1, and similarly the equation<br />

sinv =<br />

Tu + S<br />

sinhu<br />

determines the path through t2. We have to be careful now to make sure we get<br />

the correct branch <strong>of</strong> arcsinv. It turns out that we must make different choices<br />

for u < u1 and u > u1. A similar formula will hold for negative u.<br />

Set u0 = S<br />

T , and ∆ = u1 − u0. It can be shown that sinv is decreasing in<br />

[u0,∞) and we cut <strong>of</strong>f the integral at u0 and use the constant value <strong>of</strong> v = v0 = π<br />

between [−u0,u0]. Furthermore, the bulk part <strong>of</strong> the integral seems to be<br />

roughly contained in the interval [u1 − ∆,u1 + ∆] (from numerical experiments).<br />

We thus have five intervals:<br />

]∞,−u1 − ∆], [−u1 − ∆,−u0], [−u0,u0], [u0,u1 + ∆], [u1 + ∆,∞[,<br />

which will be treated separately. In the following sections arcsin always de-<br />

note the principal branch <strong>of</strong> arcsine, i.e. arcsinx ∈ [− π π<br />

2 , 2 ].<br />

In all <strong>of</strong> the intervals we need to find expressions in terms <strong>of</strong> u for the<br />

functions v, dv<br />

du , φ(u) and cosh(u + iv(u)).<br />

The interval u > u1<br />

Here we have<br />

and<br />

dv<br />

du<br />

v = arcsin<br />

<br />

Tu − S<br />

,<br />

sinhu<br />

T sinhu − (Tu − S)coshu<br />

= <br />

sinhu sinh 2 u − (Tu − S) 2<br />

.<br />

It is also clear that cosh(u + iv) = coshucosv + isinhusinv = coshucosv +<br />

i[Tu − S], hence<br />

The interval u < u1<br />

φ(u + iv) = −xcoshucosv − Rv − i[sinhusinv − Ru]<br />

= −xcoshucosv − Rv − i[x(Tu − S) − Ru]<br />

= −xcoshucosv − Rv + ixS.<br />

Note that v1 = π<br />

2 so in this interval we have to change the branch <strong>of</strong> arcsine,<br />

i.e.<br />

<br />

Tu − S<br />

v = π − arcsin ,<br />

sinhu<br />

144


and<br />

dv<br />

du<br />

T sinhu − (Tu − S)coshu<br />

= −<br />

sinhu sinh 2 .<br />

u − (Tu − S) 2<br />

As above, we also have cosh(u + iv) = coshucosv + i[Tu − S], and<br />

The interval −u0 ≤ u ≤ u0<br />

φ(u) = −xcoshucosv − Rv + ixS.<br />

In this interval we keep v = v0 = π constant, hence<br />

and<br />

Here cosh(u + iv) = −coshu, and<br />

v = π,<br />

dv<br />

= 0.<br />

du<br />

φ(u) = −xcoshucosπ − Rπ + iRu = −Rπ + xcoshu + iRu.<br />

The interval −u1 ≤ u ≤ u0<br />

Since we have π<br />

2 ≤ v ≤ π here also we have to stick to the same branch <strong>of</strong><br />

arcsine, hence<br />

<br />

Tu + S<br />

v = π − arcsin ,<br />

sinhu<br />

and<br />

dv<br />

du<br />

T sinhu − (Tu + S)coshu<br />

=<br />

sinhu sinh 2 .<br />

u − (Tu + S) 2<br />

We get cosh(u + iv) = coshucosv + i[Tu + S], and<br />

φ(u+iv) = −xcoshucosv−Rv−i[x(Tu+S)−Ru] = −xcoshucosv−Rv−ixS.<br />

The interval u ≤ −u1<br />

Now v ≤ π<br />

2 and we can get back to the principal branch:<br />

and<br />

dv<br />

du<br />

v = arcsin<br />

<br />

Tu + S<br />

,<br />

sinhu<br />

T sinhu − (Tu + S)coshu<br />

= −<br />

sinhu sinh 2 .<br />

u − (Tu + S) 2<br />

145


We have cosh(u + iv) = coshucosv + i[Tu + S], and<br />

T ≤ 1 :<br />

φ(u + iv) = −xcoshucosv − Rv − i[xsinhusinv − Ru]<br />

= −xcoshucosv − Rv − i[x(Tu + S) − Ru]<br />

= −xcoshucosv − Rv − ixS.<br />

If T ≤ 1 we define T = R<br />

x = sinα, and we can see that φ ′ (t) = 0⇔ t = iα, that<br />

is, the stationary point is on the imaginary axis. We also see that φ ′′ (iα) =<br />

−xcosh(iα) = −xcosα = −x √ 1 − T 2 = − √ x 2 − R 2 ≤ 0. Since the point t =<br />

iα is a maximum (i.e. φ ′′ (iα) < 0) we see that there might be trouble as<br />

T → 1; numerically it turns out that a safe restriction is<br />

|R − x| > 5 2 3 R 1 3 ,<br />

and for |R − x| < 5 2 3 R 1 3 we will use the path in the next section instead. The<br />

path <strong>of</strong> stationary phase that passes through iα is given by ℑφ(u + iv) =<br />

ℑφ(iα) ⇔ Ru − xsinhusinv = 0, that is<br />

and<br />

dv<br />

du<br />

sinv = Ru Tu<br />

=<br />

xsinh sinhu ,<br />

sinhu − ucoshu<br />

= ±R<br />

sinhu (xsinhu − Ru)(xsinhu + Ru) ,<br />

where +is for u > 0 and − for u < 0. We also have<br />

and<br />

The case T ≈ 1 :<br />

cosh(u + iv) = coshucosv + iTu,<br />

φ(u) = iR(u + iv) − xcoshucosv<br />

= −xcoshucosv − Rv + i[−xsinhusinv + Ru]<br />

= −xcoshucosv.<br />

This case is considered separately because <strong>of</strong> the numerical reasons hinted at<br />

above, and we will treat T close to 1 as if exactly 1. If we didn’t do that we<br />

might for example end up with sinv > 1. The test used to indicate that T is<br />

close to 1 is that<br />

|x − R| < min( R<br />

3 ,52 3 R 1 3 ).<br />

146


We get<br />

and as before<br />

and<br />

sinv = u<br />

sinhu ,<br />

dv<br />

du =<br />

sinhu − ucoshu<br />

sinhu (sinhu − u)(sinhu + u) ,<br />

φ(u) = −xcoshucosv − Rv + i[−xsinhusinv + Ru]<br />

= −xcoshucosv − Rv + iu[R − x].<br />

We also get cosh(u + iv) = coshucosv + iu.<br />

4.4 The integrand<br />

When we review the previous sections, we see directly that in all cases we<br />

have v(−u) = v(u), dv dv<br />

du (−u) = du (u) and φ(−u) = φ(−u). Using<br />

cosh(−u + iv) = cosh(u + iv),<br />

and the fact that we see that Ψ is analytic in its third argument we have<br />

<br />

ψ(−u) =Ψ −l, 1<br />

<br />

,x(1 + cosh(−u + iv(−u)) = Ψ −l,<br />

2 1<br />

<br />

,x(1 + cosh(u + iv)<br />

2<br />

=ψ(u).<br />

Note that when using this formula we must be careful since Ψ is also in general<br />

multi-valued in its third argument (e.g. Ψ − 1 1<br />

2 , 2 ,z = √ z). It is customary to<br />

use the principal branch, i.e. with a branch-cut along the negative real-axis,<br />

but in the case T ≥ 1 however, the argument x(1 + cosh(−u + iv(−u)) passes<br />

through the negative real axis at u = ±u0 and we must make sure we get the<br />

correct branch choice.<br />

We thus have the following reflectional symmetries:<br />

v(−u) = v(u),<br />

dv<br />

(−u)<br />

du<br />

=<br />

dv<br />

du ,<br />

φ(−u) = φ(u),<br />

ψ(−u) = ψ(u),<br />

and using these we can rewrite the complex integral (4.1) as a real integral. To<br />

147


simplify the following equations we introduce the notation<br />

where<br />

We now have<br />

I = 1<br />

2<br />

= 1<br />

2<br />

= 1<br />

2<br />

= 1<br />

2<br />

=<br />

=<br />

∞<br />

−∞<br />

∞<br />

−∞<br />

0<br />

−∞<br />

∞<br />

0<br />

∞<br />

0<br />

∞<br />

0<br />

e φ(u) ψ(u) = P(u) = R(u)e iΘ(u) ,<br />

Θ(u) = ℑφ(u) + Argψ(u), and<br />

R(u) = |ψ(u)|e ℜφ(u) .<br />

e φ(t) ψ(t)dt<br />

<br />

P(u) 1 + i dv<br />

<br />

du<br />

du<br />

<br />

P(u) 1 + i dv<br />

<br />

du +<br />

du<br />

1<br />

<br />

∞<br />

P(u) 1 + i<br />

2 0<br />

dv<br />

<br />

du<br />

du<br />

<br />

P(−u) 1 − i dv<br />

<br />

+ P(u) 1 + i<br />

du<br />

dv<br />

<br />

du<br />

du<br />

1<br />

i<br />

dv<br />

(P(u) + P(−u)) + (P(u) − P(−u))<br />

2 2 du du<br />

<br />

R(u) cos(Θ(u)) − sin(Θ(u)) dv<br />

<br />

du,<br />

du<br />

which is now a completely real integral. When we split it up for the different<br />

cases we get explicit expressions for the various factors.<br />

T ≥ 1 :<br />

Here we have<br />

where now<br />

and<br />

148<br />

I1 =<br />

u0<br />

0<br />

= e −Rπ<br />

I =<br />

u0<br />

0<br />

= I1 + I2,<br />

∞<br />

+<br />

u0<br />

<br />

R(u) cos(Θ(u)) − sin(Θ(u)) dv<br />

<br />

du<br />

du<br />

u0<br />

0<br />

e xcoshu |ψ(u)|{cos(Ru + Argψ(u))}du,<br />

<br />

∞<br />

I2 = R(u) cos(Θ(u)) − sin(Θ(u))<br />

u0<br />

dv<br />

<br />

du,<br />

du


where <br />

Θ(u) = xS + Argψ(u),<br />

T ≤ 1 :<br />

Here<br />

T ≈ 1 :<br />

I =<br />

=<br />

=<br />

−∞<br />

∞<br />

∞<br />

−∞<br />

∞<br />

0<br />

∞<br />

0<br />

<br />

P(u)<br />

R(u)<br />

e φ(u) |ψ(u)|<br />

R(u) = |ψ(u)|e −xcoshucosv−Rv .<br />

1 + i dv<br />

du<br />

The relevant integral is given by<br />

<br />

∞<br />

I = P(u) 1 + i dv<br />

<br />

du =<br />

du<br />

=<br />

=<br />

0<br />

∞<br />

0<br />

R(u)<br />

e φ(u) |ψ(u)|<br />

<br />

du =<br />

<br />

cos(Θ(u)) − sin(Θ(u)) dv<br />

du<br />

<br />

du.<br />

<br />

cos(Argψ(u)) − sin(Argψ(u)) dv<br />

<br />

du.<br />

du<br />

<br />

cos(Θ(u)) − sin(Θ(u)) dv<br />

du<br />

<br />

du.<br />

<br />

cos(uδ + Argψ(u)) − sin(uδ + Argψ(u)) dv<br />

<br />

du.<br />

du<br />

These integrals we have obtained can now be dealt with using standard<br />

methods for numerical integration, e.g. Gaussian quadrature once we have<br />

satisfying algorithms to compute the function ψ.<br />

T ≥ 1 :<br />

Here we have<br />

where now<br />

I1 =<br />

u0<br />

0<br />

= e −Rπ<br />

I =<br />

u0<br />

0<br />

= I1 + I2,<br />

∞<br />

+<br />

u0<br />

<br />

R(u) cos(Θ(u)) − sin(Θ(u)) dv<br />

<br />

du<br />

du<br />

u0<br />

0<br />

e xcoshu |ψ(u)|{cos(Ru + Argψ(u))}du,<br />

149


and<br />

<br />

∞<br />

I2 = R(u) cos(Θ(u)) − sin(Θ(u))<br />

u0<br />

dv<br />

<br />

du,<br />

du<br />

where <br />

Θ(u) = xS + Argψ(u),<br />

T ≤ 1 :<br />

Here<br />

T ≈ 1 :<br />

I =<br />

=<br />

=<br />

−∞<br />

∞<br />

∞<br />

−∞<br />

∞<br />

0<br />

∞<br />

0<br />

<br />

P(u)<br />

R(u)<br />

e φ(u) |ψ(u)|<br />

R(u) = |ψ(u)|e −xcoshucosv−Rv .<br />

1 + i dv<br />

du<br />

The relevant integral is given by<br />

<br />

∞<br />

I = P(u) 1 + i dv<br />

<br />

du =<br />

du<br />

=<br />

=<br />

0<br />

∞<br />

0<br />

R(u)<br />

e φ(u) |ψ(u)|<br />

<br />

du =<br />

<br />

cos(Θ(u)) − sin(Θ(u)) dv<br />

du<br />

<br />

du.<br />

<br />

cos(Argψ(u)) − sin(Argψ(u)) dv<br />

<br />

du.<br />

du<br />

<br />

cos(Θ(u)) − sin(Θ(u)) dv<br />

du<br />

<br />

du.<br />

<br />

cos(uδ + Argψ(u)) − sin(uδ + Argψ(u)) dv<br />

<br />

du.<br />

du<br />

These integrals we have obtained can now be dealt with using standard<br />

methods for numerical integration, e.g. Gaussian quadrature once we have<br />

satisfying algorithms to compute the function ψ.<br />

4.5 Evaluation <strong>of</strong> the Integrand<br />

As a final step to be able to evaluate the integrals in the previous section we<br />

need efficient algorithms for the function Ψ −l; 1<br />

2 ,z where l is real and z is<br />

a complex number. In most papers on numerical algorithms for Ψ(a,c,z) the<br />

values <strong>of</strong> z are restricted to the reals, cf. e.g. [96]. Algorithms without this<br />

restriction can be found in e.g. [102].<br />

150


In this chapter we will outline the three different approaches which seems<br />

to yield the most efficient algorithms in different z-domains. The methods<br />

include power-series for small |z|, backward-recursion for middle range |z|<br />

and an asymptotic series expansion for large |z|. Each <strong>of</strong> them is evaluated to<br />

at least 12-14 significant digits in the relevant interval.<br />

It should also be noted that for the special case when l = 1<br />

2 (which will be<br />

used later) we have <br />

Ψ − 1 1<br />

,<br />

2 2 ,z<br />

<br />

= √ z,<br />

and another special case is (<strong>of</strong> course) weight 0, that is l = 0, where we have<br />

for all c<br />

Ψ[0,c,z] = Φ[0,c,z] = 1.<br />

These special cases can be easily deduced from the power series expression in<br />

the next subsection.<br />

4.5.1 The power series<br />

It is known that we can express Ψ as a sum <strong>of</strong> two Φ functions (cf. [99, p. 20]<br />

or [29, p. 257]):<br />

Ψ(a,c,z) =<br />

Γ(1 − c)<br />

Γ(c − 1)<br />

Φ(a,c,z) + z<br />

Γ(a − c + 1) Γ(a)<br />

1−c Φ(a − c + 1,2 − c,z),<br />

hence we actually only need an algorithm for Φ (the algorithm for Γ, using<br />

Stirling’s formula is standard). We know that Φ has a power-series expansion<br />

valid for all a, z and c not a negative integer (if c is a negative integer it is valid<br />

if a = c)<br />

Φ(a,c,z) =<br />

∞<br />

∑<br />

n=0<br />

(a)n<br />

where (a)n is the Pochammer symbol, defined as (a)k = a(a+1)···(a+k−1),<br />

for a complex (not a negative integer > −k +1) and k a positive integer. It can<br />

also be written as (a)k = Γ(a+k)<br />

Γ(a) . We see that for non-exceptional a and c (i.e.<br />

c not very close to a negative integer) the power series should converge in a<br />

similar manner as the power series for the exponential function. A few special<br />

(c)n<br />

z n<br />

n! ,<br />

151


cases worth mentioning are<br />

Φ(0,c,z) = 1,<br />

Φ(a,a,z) = e z ,<br />

Φ(−1, 1<br />

,z)<br />

<br />

2<br />

Φ 1,<br />

= 1 − 2z,<br />

1<br />

2 ,z<br />

<br />

= √ π<br />

∞<br />

∑<br />

n=0<br />

z n<br />

Γ 1<br />

2 + n.<br />

To evaluate the power series one can use for example a direct formula, recursion,<br />

or a continued-fraction method (c.f. [71, p. 185]). We set<br />

A−1 = A0 = 1, A1 = 1 + az<br />

c ,<br />

ri =<br />

(a + i − 1)<br />

, i > 1,<br />

i(c + i − 1)<br />

Ai = Ai−1 + (Ai−1 − Ai−2)riz, for i > 1.<br />

Then the i − thapproximation to Φ(a,c,z) is given by Ai, and we quit the<br />

summation whenever |Ai − Ai−1| < ε.<br />

The special case <strong>of</strong> c = 1<br />

2<br />

Since what we in general want to evaluate is Ψ(a, 1<br />

2 ,z) we can formulate the<br />

specific formula for this case,<br />

<br />

Ψ a, 1<br />

2 ,z<br />

<br />

Γ<br />

=<br />

<br />

1<br />

2<br />

Γ a + 1<br />

<br />

Φ a,<br />

2<br />

1<br />

2 ,z<br />

<br />

+ Γ− 1<br />

<br />

2 √<br />

zΦ a +<br />

Γ(a)<br />

1 3<br />

,<br />

2 2 ,z<br />

<br />

Γ<br />

=<br />

<br />

1<br />

2<br />

Γ a + 1<br />

∞ (a)n<br />

∑ <br />

1<br />

2 n=0 2 n n!zn + Γ− 1<br />

<br />

2 √ ∞ 1 a + 2<br />

z<br />

Γ(a) ∑ n<br />

zn<br />

3<br />

n=0 2 n! n<br />

Γ<br />

=<br />

<br />

1<br />

2<br />

Γ a + 1<br />

∞ Γ(a + n)Γ<br />

∑<br />

2 n=0<br />

<br />

1<br />

2<br />

Γ(a)Γ 1<br />

2 + nΓ(n + 1) zn<br />

+ Γ− 1<br />

<br />

2 √ ∞ Γ<br />

z<br />

Γ(a) ∑<br />

n=0<br />

a + 1<br />

2 + nΓ <br />

3<br />

2<br />

Γ a + 1<br />

<br />

3<br />

2 Γ 2 + nΓ(n + 1) zn<br />

Γ<br />

=<br />

<br />

1 2<br />

2<br />

Γ a + 1<br />

<br />

∞<br />

∑ Γ(a + n) −<br />

2 Γ(a) n=0<br />

√ z Γa + 1<br />

<br />

2 + n<br />

×<br />

1<br />

2 + n<br />

×<br />

z n<br />

Γ(n + 1)Γ n + 1<br />

.<br />

2<br />

This sum can now be evaluated by any standard method. Define<br />

152<br />

an =<br />

Γ(a + n)<br />

Γ(n + 1)Γ n + 1<br />

, bn =<br />

2<br />

Γ a + 1<br />

2<br />

Γ(n + 1)Γ n + 1<br />

2<br />

+ n<br />

<br />

1<br />

2 + n,


and compute the sum<br />

∞ <br />

∑ an −<br />

n=0<br />

√ n<br />

zbn z .<br />

Since the an and bn are independent <strong>of</strong> z and this sum is evaluated many times<br />

it is preferable to make a table <strong>of</strong> an and bn. These can either be computed<br />

by recursion to get good speed, or since we need only tabulate them once by<br />

direct evaluation.<br />

Reflection<br />

The power series is obviously valid in a circular region in the plane, but experimentally<br />

it seems more efficient to use Kummer’s transformation for values<br />

<strong>of</strong> z in the left half-plane:<br />

Φ(a,c,z) = e z Φ(c − a,c,−z),<br />

which gives the following formula for Ψ:<br />

Ψ(a,c,−z) =<br />

e −z<br />

<br />

Γ(1 − c)<br />

Γ(c − 1)<br />

Φ(c − a,c,z) + (−z)1−c Φ(1 − a,2 − c,z) ,<br />

Γ(a − c + 1) Γ(a)<br />

where Φ can be evaluated as above.<br />

4.5.2 Recursion<br />

It is well-known that both confluent hypergeometric functions satisfy many<br />

recursion relations (cf. [29, pp. 254, 257-258]), and one can use these to<br />

formulate recurrence equations, cf. e.g. [96, pp. 64-65]. If we set un =<br />

(a)nΨ(a + n,c,z) we get the following equation:<br />

(a + n − 1)un−1 + (c − 2a − z − 2n)un + (a + n + 1 − c)un+1 = 0, (4.4)<br />

which is valid for all values <strong>of</strong> a,c,z and for all integers n ≥ 1. What we want<br />

to compute is now u0, the first term in the recursion, hence we can’t use regular<br />

forward-recursion. What we have to do is basically to start with some values<br />

<strong>of</strong> uN and uN−1 and then go backwards. The standard way to do this is by a so<br />

called Miller algorithm (cf. e.g. [32, 73, 74] or in this particular case [96] and<br />

[102]). In order to use this kind <strong>of</strong> algorithm we need a normalization, and in<br />

this case we use<br />

S =<br />

∞<br />

∑ λnun = z<br />

n=0<br />

−a ,<br />

153


where λn = (−1) n<br />

<br />

c − a − 1<br />

n<br />

<br />

= (−1) n Γ(c−a)<br />

Γ(c−a−n)n! . This follows from the<br />

binomial theorem for general exponents. What we do next is to rewrite the<br />

recursion equation (4.4) as<br />

where<br />

an =<br />

un+1 + anun + bnun−1 = 0,<br />

c − 2a − z − 2n<br />

1 + n + a − c<br />

and bn =<br />

a + n − 1<br />

1 + n + a − c .<br />

There are several choices in the implementation and the one we use is basically<br />

[32, p. 37]. This is actually a continued-fraction method for the quotients<br />

un+1/un. For a fixed positive integer N > 0 we set<br />

and then recursively set<br />

s (N)<br />

N<br />

r (N)<br />

n−1 =<br />

= r(N)<br />

N<br />

−bn<br />

an + rn<br />

= 0,<br />

, and<br />

s (N)<br />

n−1 = rn−1(λn + s (N)<br />

n ),<br />

for n = N,N − 1,N − 2,...1. In the final step we get an approximation to u0:<br />

u (N)<br />

0 =<br />

S<br />

λ0 + s (N)<br />

0<br />

To get an approximation to u0 with an error smaller than ε we choose a starting<br />

integer M and repeat the process for the sequence <strong>of</strong> N = M,M +1,... and stop<br />

when the difference |u (N)<br />

0 − u(N+1)<br />

0 | < ε. To avoid unnecessary computations<br />

the choice <strong>of</strong> M is important, and should be tabulated in advance.<br />

Note that λn satisfies a recursion relation: λ0 = 1, and λn = − c−a−n<br />

n λn−1,<br />

and if we compute {λ j} M 1 in advance we only need to compute an additional<br />

λN in each step.<br />

From the formula for an and bn we see that the only points we should avoid<br />

are those where 1 + a − c is a negative integer. Hence this method should be<br />

used for the sake <strong>of</strong> precision when c is close to an integer but 1 + a − c is not<br />

a negative integer.<br />

154<br />

.


4.5.3 Asymptotic series<br />

The following asymptotic expansion for Ψ is well-known (see [29, p. 278])<br />

Ψ(a,c,z) =z −a<br />

<br />

N<br />

∑ (−1)<br />

n=0<br />

n (a)n(a − c + 1)n<br />

n!zn + O(z −N−1 <br />

) , as |z| → ∞,<br />

for |Argz| ≤ 3π<br />

2 ,<br />

and when evaluated by forward-recursion on the coefficients the convergence<br />

is fast for large |z|. As is usual when evaluating asymptotic series we sum<br />

the series up to the term N = N0 where the terms starts to increase in absolute<br />

value.<br />

4.5.4 Comparison<br />

We made a spot-check <strong>of</strong> the methods for computing Ψ(−l, 1<br />

2<br />

,x + iy) over<br />

a grid where |l| ≤ 1 and 1 ≤ x ≤ 100 and for each (l,x) we let y vary between<br />

1 and lnx to simulate the behaviour <strong>of</strong> cosh(x + iy) at the paths above,<br />

where basically the imaginary part is sinvsinhu = O(u) and the real part is<br />

cosvcoshu = O(eu ). In this region we first <strong>of</strong> all compared the three methods<br />

with each other with respect to speed and we also made spot-check against<br />

Maple and PARI (with internal precision set to 100 decimal places). We found<br />

that the program was accurate up to 12-13 significant digits in the entire region<br />

when we choose the algorithms as follows:<br />

0 ≤ x ≤ 25 Power series<br />

25 ≤ x ≤ 35 Recursion<br />

35 ≤ x Asymptotic series<br />

It should be noted that the recursive algorithm actually provides 12-13 decimal<br />

places (depending on the termination condition <strong>of</strong> course) for x ≥ 0.1, but is<br />

in general very slow for small x. By further sporadic spot-checking we also<br />

drew the conclusion that the Asymptotic series is accurate up to at least 12-<br />

13 decimal places when |z| ≥ 35. Negative values <strong>of</strong> x was checked against<br />

positive values by the reflection relation stated above that relates Ψ(a,c,−x)<br />

with Ψ(a,c,x).<br />

4.6 Evaluation <strong>of</strong> the Integrals<br />

The infinite integrals I were first truncated to finite integrals using the fact that<br />

,x) grows at most like a polynomial in x <strong>of</strong> degree ≤ |l| and that eℜφ<br />

Ψ(−l, 1<br />

2<br />

155


is decreasing like e−ex. To obtain the bound on Ψ(−l, 1<br />

2 ,x) study the integral<br />

representation<br />

Ψ(−l, 1<br />

<br />

1 ∞<br />

,x) = e<br />

2 Γ(−l) 0<br />

−xt t −l−1 (1 +t) l− 1 2 dt,<br />

do a variable substitution u = xt and compare with the integral representation<br />

and known bounds for the Γ function.<br />

The actual bounds we use are verified by the numerical experiments. The<br />

finite interval is then divided up in to smaller intervals (according to the value<br />

<strong>of</strong> R) and on each <strong>of</strong> the smaller intervals we evaluate the integral using a<br />

16-point Gaussian quadrature rule.<br />

4.7 Remarks on the K-Bessel function<br />

We know that the special case l = 0 in Wl,iR(x) gives us essentially the K-<br />

Bessel function, KiR(x) (i.e. W0,iR(2x) =<br />

2x<br />

π KiR(x)). The algorithm de-<br />

scribed in this paper is then identical to Hejhal’s algorithm ([42]), generalized<br />

by Avelin to Ks(x) for complex s ([9]).<br />

Note that<br />

KiR(x) =<br />

<br />

π<br />

2x e−x (2x) 1 2 +iR <br />

1<br />

Ψ + iR,1 + 2iR;2x ,<br />

2<br />

but even though the arguments to the Ψ−function are complex it turns out that<br />

in this special case most <strong>of</strong> the algorithms we used to compute the Ψ-function<br />

with real arguments can be trivially modified and actually show good convergence<br />

properties. Hence we can use both the backwards recursion (Millertype)<br />

algorithm and the power series to compute the K-Bessel function.<br />

One particularly nice power series method for K-Bessel function is the following.<br />

From the expression <strong>of</strong> KiR as a sum <strong>of</strong> IrR and I−iR it is easy to see the<br />

power series defining KiR from the known power series for IiR. From [33, pp.<br />

399-400], and breaking out the dominating factors from the expansion we see<br />

that<br />

with<br />

156<br />

1 π<br />

KiR(x) =<br />

|Γ(1 + iR)| 2sinh(πR)<br />

fk = −i|Γ(1 + iR)|<br />

∞<br />

∑ ck fk<br />

k=0<br />

<br />

x −iR<br />

2<br />

Γ(k + 1 − iR) −<br />

<br />

x iR<br />

<br />

2<br />

Γ(k + 1 + iR)


and<br />

and<br />

rk = R|Γ(1 + iR)|<br />

<br />

x −iR<br />

2<br />

Γ(k + 1 − iR) +<br />

<br />

x iR<br />

<br />

2<br />

Γ(k + 1 + iR)<br />

,<br />

ck = 1<br />

k!<br />

x 2<br />

4<br />

k<br />

.<br />

The recursive relations are not affected by this scaling and are still<br />

fk = k fk−1 + rk−1<br />

k2 + R2 rk =<br />

1<br />

k2 + R2 [(2k − 1)rk−1 − rk−2]<br />

ck = x2 ck−1<br />

4 k<br />

for k ≥ 1 where the initial data now are<br />

f0 = −2sinθ,<br />

r0 = 2Rcosθ, and<br />

r1 =<br />

R<br />

(cosθ + Rsinθ)<br />

1 + R2 with θ = Rln( x<br />

2 ) − ArgΓ(1 + iR). Now we can just compute the sum by standard<br />

forward-recursion and multiply with the Γ and sinh factors at the end.<br />

(And for large R avoid cancellation and division by infinity by simply take<br />

both factors into account simultaneously). (See also [98, pp. 377-379].)<br />

4.8 Numerical verification<br />

The numerical verification <strong>of</strong> the algorithm has been done by comparing the<br />

values <strong>of</strong> Wl,iR(x) against Maple with 25 digits precision. The comparison has<br />

been done over 2 separate grids in the (R,x,l) space:<br />

1. 200000 pseudo random points (Ri,xi,li) in the cube<br />

[0,20] × [0,20] × [0,6], and<br />

2. 100000 pseudo random points (Ri,xi,li) in the cube<br />

[20,100] × [0,20] × [0,6].<br />

The reason for the choice <strong>of</strong> [0,6] as the l-interval is because that is the range<br />

we will use in the application <strong>of</strong> this algorithm to the problem <strong>of</strong> computing<br />

<strong>Maass</strong> <strong>waveforms</strong> with nonzero weight.<br />

157


We consider the relative errors<br />

εi = Wli,iRi (xi) Fortran −Wli,iRi (xi) Maple <br />

Wli,iRi (xi) Maple −1<br />

as corresponding to a set <strong>of</strong> (pseudo)random variables. As such we consider<br />

their mean and variance, apart from the obvious study <strong>of</strong> their absolute maximum.<br />

We also consider the mean and standard deviance <strong>of</strong> the absolute values<br />

|εi|.<br />

1. In the first grid we have<br />

2. In the second grid we have<br />

max{|εi|} = 5.3E − 10,<br />

mean{εi} = −2.0E − 15,<br />

mean{|εi|} = 4.2E − 14,<br />

std{|εi|} = 3.5E − 12,<br />

std{εi} = 3.5E − 12.<br />

max{|εi|} = 1.7E − 09,<br />

mean{εi} = 2.2E − 14,<br />

mean{|εi|} = 8.1E − 14,<br />

std{|εi|} = 5.5E − 12,<br />

std{εi} = 5.5E − 12.<br />

In connection with this it might also be enlightening to know that CPU time<br />

required for Maple to make the computations over the first grid is approximately<br />

27 hours, where my Fortran program only need about 15 minutes (on<br />

a 1533MHz CPU). Over the second grid the needed CPU times are approximately<br />

25h45m for Maple and 33m for Fortran (on a 1533MHz CPU). These<br />

figures are not exact since there were other programs running simultaneously<br />

on this computer too, but they serve well as a hint <strong>of</strong> the effectiveness <strong>of</strong> the<br />

algorithm.<br />

158<br />

,


Swedish Summary<br />

<strong>Maass</strong>vågfunktioner ur ett beräkningsperspektiv<br />

Detta är en avhandling i experimentell matematik, och vi befinner oss således<br />

i en situation där den “rena” matematiken drar nytta av numeriska resultat för<br />

att bygga hypoteser.<br />

En kortfattad beskrivning av vad som kännetecknar experimentell matematik<br />

är att man med hjälp av teoretiska resonemang konstruerar numeriska<br />

modeller (algoritmer) vilka sedan används för att generera information. Utifrån<br />

de numeriska resultaten bygger man sedan upp hypoteser och förmodanden<br />

(för en mer detaljerad beskrivning se [16]).<br />

Förhoppningen är naturligtvis att hypoteserna sedan skall kunna bevisas<br />

strikt matematiskt, men även i de fall då det nuvarande matematiska maskineriet<br />

inte förmår detta kan de numeriska resultaten ge värdefulla insikter<br />

och inspirera till fortsatta studier.<br />

Avhandlingens tre huvudkapitel är alla utformade på samma sätt; merparten<br />

ägnas åt att presentera den teori som sedan används för att konstruera algoritmer.<br />

Efter en beskrivning av hur algoritmerna kan utformas presenteras de<br />

experimentella resultat vi kommit fram till.<br />

Det problem vi studerat är att konstruera de matematiska objekt som kallas<br />

“<strong>Maass</strong>vågfunktioner” (efter Hans <strong>Maass</strong> (1911-1992)). Matematiska tillämpningar<br />

av dessa vågfunktioner finns främst inom analytisk talteori. Det finns<br />

även fysikaliska kopplingar, främst då till matematiska modeller av kvantkaos<br />

([35, 14]), men även till kosmologi ([8]).<br />

Något om den fysikaliska bakgrunden<br />

Betrakta en (kvantmekanisk) partikel med massa m0 som rör sig fritt på en<br />

yta, S . Enligt kvantmekaniken beskrivs detta system av en vågfunktion, ψ,<br />

vilken talar om för oss hur stor sannolikheten är att påträffa partikeln i ett visst<br />

område. Vågfunktionens utbredning, för en partikel med energin E = ¯hv, ges<br />

i sin tur av Schrödingerekvationen:<br />

i¯h ∂ψ<br />

∂t<br />

¯h2<br />

= − ∆ψ = Eψ,<br />

2m0<br />

159


där ∆ är den så kallade Laplace-Beltramioperatorn på ytan. (T.ex. ges den i<br />

planet av ∆ψ = ψ ′′<br />

xx + ψ ′′<br />

yy). Vi antar nu att det går att skriva ψ som en produkt<br />

av en tidsberoende del, e−ivt , och en tidsoberoende del, φ(P), P ∈ S . Vi<br />

får då en ekvation för den tidsoberoende delen (den tidsoberoende Schrödingerekvationen,<br />

eller Laplaces ekvation):<br />

∆φ + λφ = 0, (⋆)<br />

där vi har satt λ = 2m0E<br />

¯h 2 (egenvärdet), och eftersom sannolikheten för att partikeln<br />

befinner sig på ytan är ett måste vi dessutom ha<br />

<br />

|φ| 2 dP = 1.<br />

S<br />

Om partikeln är laddad och vi kopplar på ett konstant magnetiskt fält, B, på<br />

ytan kommer den tidsoberoende delen av vågfunktionen att få ett “tvist” (spin)<br />

och φ måste nu uppfylla en variant av Laplaces ekvation:<br />

∆Bφ + λφ = 0, (⋆⋆)<br />

där ∆B är en modifierad Laplaceoperator, (i vårt fall blir den ∆B = ∆ − iyB ∂<br />

∂x<br />

med B proportionell mot den magnetiska fältstyrkan).<br />

Observera att vågekvationen, (⋆) är “universell”, i den meningen att den<br />

beskriver många olika former av vågrörelse. T.ex. om S får beskriva ett<br />

trumskinn, så kommer (⋆) att beskriva utseendet hos en stående våg på trumskinnet,<br />

med en frekvens proportionell mot √ λ. Bastoner svarar alltså mot<br />

små värden på λ och ljusare toner mot högre värden på λ.<br />

I den här avhandlingen betraktar vi ekvationerna (⋆) och (⋆⋆) för vissa<br />

speciella ytor S , nämligen sådana med konstant negativ krökning (vi säger<br />

att ytan är hyperbolisk), ändlig area samt ett ändligt (positivt) antal “punkteringar”<br />

eller “spetsar” (dvs vi har plockat bort ett ändligt antal punkter från<br />

ytan och “dragit ut dem” till oändligheten). Dessutom måste ytorna vara<br />

speciella på ett “aritmetiskt” sätt, vilket i det här fallet betyder att de ges av en<br />

delgrupp till den modulära gruppen. Vi tillåter också ytor med en sorts “hörn”.<br />

Lösningarna till (⋆) och (⋆⋆) på sådana ytor kallas för <strong>Maass</strong>vågfunktioner.<br />

En orsak till att hyperboliska ytor, är intressanta i detta sammanhang är<br />

att den klassiska dynamiken för en fri partikel på sådana är (starkt) kaotisk.<br />

Genom att studera lösningarna till (⋆) på en sådan yta kan vi sålunda komma<br />

till insikt om det finns “spår” av detta klassiska kaos kvarlämnat i det kvantiserade<br />

systemet. Denna typ av undersökningar är en del av det forskningsområde<br />

som kallas kvantkaos. Då vågfunktionerna i sig inte är kaotiska visar sig<br />

det klassiska kaoset i bl.a. energinivåernas (egenvärdenas) fördelning samt<br />

utseendet hos de vågfunktioner φ som svarar mot höga energier. I Figur 2<br />

ser du ett exempel på en <strong>Maass</strong>vågfunktion, ϕ med ett egenvärde λ ≈ 10000.<br />

160


(a) visar nodlinjerna (dvs. de punkter där ϕ(z) = 0), och i (b) kan du se sannolikhetsfunktionen<br />

som associeras till ϕ, dvs |ϕ(z)| 2 . De ljusa områdena i<br />

(b) visar de platser i planet där den kvantmekaniska partikeln gärna hänger<br />

omkring.<br />

Det huvudsakliga syftet med arbetet som ligger till grund för denna avhandling<br />

var att ta fram stabila algoritmer för att beräkna egenvärden, λ, samt<br />

motsvarande vågfunktioner, φ, till ovanstående problem.<br />

Kort om Hyperbolisk Geometri<br />

För att åskådliggöra negativ krökning (hyperbolisk geometri) brukar man som<br />

matematiker arbeta i t.ex. övrehalvplansmodellen, dvs<br />

H = {x + iy ∈ C|y > 0},<br />

och där införa ett avstånds- samt areabegrepp med hjälp av bågelementet<br />

och areaelementet<br />

ds = |dz|<br />

y ,<br />

dµ = dxdy<br />

.<br />

y2 Den kortaste vägen mellan två punkter z1,z2 ∈ H får man genom att följa<br />

den geodet som passerar genom dessa. I vanliga fall (dvs då vi inte har någon<br />

krökning) så är geodeten den unika räta linje som går mellan punkterna. I det<br />

hyperboliska övre halvplanet ges geodeterna istället av halvcirklar samt räta<br />

linjer (vilka kan betraktas som generaliserade halvcirklar), vilka står vinkelrätt<br />

mot x-axeln. Se Figur S1 för en illustration av geodeterna γ1 och γ2 som<br />

Figure S1: Övre halvplanet<br />

H<br />

γ3<br />

y<br />

z4<br />

z3<br />

γ1<br />

γ2<br />

z2<br />

z1<br />

(a) Geodeter och en horocykel<br />

x<br />

S<br />

F1 F4 F7<br />

E<br />

e2<br />

e3 e ′ 3<br />

F6<br />

F3<br />

F5<br />

F2<br />

(b) Fundamentalområde<br />

förbinder z1 med z2 respektive z3 med z4. I figuren ger vi även ett exempel på<br />

161


en så kallad horocykel, γ3.<br />

En konsekvens av detta sätt att mäta på i det övre halvplanet är t.ex. att man<br />

inte kan ta sig ända ner till den reella linjen (dvs avståndet från varje punkt i<br />

H till reella linjen är oändligt).<br />

Ett något mer fysikaliskt sätt att tänka sig detta hyperboliska rummet på är<br />

följande (fritt efter Poincaré [82]):<br />

"Tänk dig att du lever i ett sfäriskt universum begränsad av ett “skal”. Antag att<br />

temperaturen i detta universum är proportionellt mot avståndet till skalet, dvs<br />

det är varmast i mitten och på skalet råder den absoluta nollpunkten. Föreställ<br />

dig nu att du rör dig med en hastighet proportionell mot temperaturen (detta<br />

är vad man fysikaliskt förväntar sig av på molekylär/atomär nivå). Dvs, ju<br />

närmare skalet man kommer, desto kallare blir det, desto saktare färdas man.<br />

Det är då ganska enkelt att inse att om man bestämmer sig för att försöka åka till<br />

“universums ände” (skalet) så kommer man aldrig fram hur länge man än åker<br />

(inom ändlig tid). Matematiskt sett kan man beräkna att r(t) = R(1 − e −ct ),<br />

där r(t) är avståndet från centrum vid tiden t, R är radien på skalet och c är en<br />

positiv konstant."<br />

En isometri i ett rum (där vi har ett avståndsbegrepp) är en avbildning som<br />

bevarar avstånd och areor. I vanliga “okrökta” fallet är translationer, rotationer<br />

och speglingar exempel på isometrier. Spegling är ett exempel på en<br />

orienteringsreverserande isometri, medan translation och rotation är orienteringsbevarande.<br />

I det övre hyperboliska övre halvplanet så kan varje orienteringsbevarande<br />

isometri skrivas som en Möbiusavbildning:<br />

T : z ↦→<br />

az + b<br />

, där a,b,c,d ∈ R, ad − bc = 1.<br />

cz + d<br />

En sådan avbildning associeras sedan till ett par av matriser:<br />

<br />

±<br />

a<br />

c<br />

b<br />

d<br />

,<br />

och gruppen av alla sådana isometrier skriver vi <strong>of</strong>tast som PSL(2,R) =gruppen<br />

av alla matriser <br />

a b<br />

c d , med a,b,c,d ∈ R, ad − bc = 1, och där vi identifierar<br />

matriserna <br />

a b −a −b<br />

c d och −c −d .<br />

Man vet nu att alla hyperboliska ytor kan beskrivas med hjälp av diskreta<br />

undergrupper till PSL(2,R), dvs vi kan skriva vår yta S som S = Γ\H , där<br />

Γ är en Fuchsisk grupp (en diskret grupp av Möbiusavbildningar). De ytor, S ,<br />

vi är intresserade av har dels “horn” (punkter utdragna till oändligheten, och<br />

vi tillåter dessutom att de har “hörn”. I den motsvarande gruppen, Γ, visar<br />

sig hornen som fixpunkter till paraboliska avbildningar, dvs de svarar mot<br />

matriser med absolutbeloppet av spåret = 2 (kom ihåg att spåret av en matris<br />

162


a b<br />

c d definieras som a + d). Hörnen är å andra sidan fixpunkter till elliptiska<br />

avbildningar (|spåret| < 2). Man brukar därför även säga att punkterna är<br />

paraboliska respektive elliptiska. (De paraboliska fixpunkterna kallas också<br />

spetsar). I PSL(2,R) finns förutom elliptiska och paraboliska avbildningar<br />

(som geometriskt svarar mot rotationer och translationer) även hyperboliska<br />

(|spåret| > 2). Då de andra typerna har unika fixpunkter har hyperboliska<br />

avbildningar i regel två fixpunkter i R ∪ {∞} (och de beskriver geometriskt en<br />

förflyttning (skalning) längs den geodet som förbinder dessa).<br />

För att illustrera ytan S = Γ\H brukar man ta hjälp av ett fundamentalområde,<br />

vilket i vårt fall, kan tas som ett sammanhängande område F ⊂ H vilket<br />

har två egenskaper: för det första får inga punkter inuti F avbildas till varandra<br />

med avbildningar i Γ, och för det andra så ska Γ-bilderna av F tesselera H<br />

(dvs ∪γ∈Γγ(F ) = H). Ett exempel på fundamentalområde är standardfundamentalområdet<br />

för den modulära gruppen, PSL(2,Z) (de matriser i PSL(2,R)<br />

som endast innehåller heltal) F1 = {z = x + iy ∈ H||z| ≥ 1, |x| ≤ 1<br />

2 }.<br />

I Figur S1(b) kan vi se F1 tillsammans med 6 stycken “kopior”. Observera<br />

att alla bilder, F j, har samma hyperboliska area. (Just den här samlingen<br />

av kopior “råkar” också vara ett fundamentalområde till en grupp, Γ7,12, som<br />

återkommer i kapitel 3). Om vi identifierar sidorna i F1 med avbildningarna<br />

S : z ↦→ z + 1 och E : z ↦→ −1<br />

z får vi (topologiskt sett) en sfär med en spets<br />

(i∞) samt två hörn (e2 samt e3 som identifieras med e ′ 3<br />

). Man kan visa att de<br />

sidparandeavbildningarna E och S genererar PSL(2,Z), och faktum är att i de<br />

fall vi betraktar går det alltid att välja fundamentalområdet så att de sidparande<br />

avbildningarna genererar gruppen.<br />

Sammanfattning av Kapitlen i Avhandlingen<br />

I kapitel ett studerar vi det problem som i övrehalvplansmodellen kan skrivas<br />

som:<br />

y 2 φ ′′<br />

xx(z) + φ ′′<br />

yy(z) + λφ(z) = 0, z = x + iy ∈ H , (1)<br />

φ (T z) = χ(T )φ(z),<br />

<br />

T ∈ Γ, z ∈ H , (2)<br />

|φ| 2 dµ = 1, (3)<br />

Γ\H<br />

där χ : Γ → S1 är en jämn Dirichlet karaktär (dvs χ(AB) = χ(A)χ(B) och<br />

χ(−1) = 1), och Γ är en Hecke-kongruensundergrupp till den modulära gruppen.<br />

För ett givet positivt heltal, N, definieras Hecke-kongruensgruppen,<br />

Γ0(N), på följande sätt:<br />

<br />

<br />

Γ0(N) =<br />

a<br />

c<br />

b<br />

d<br />

|a,b,c,d ∈ Z, ad − bc = 1, c ≡ 0 mod N .<br />

163


Vi beskriver kortfattat algoritmen. Givet ett egenvärde, λ = 1<br />

4 + R2 , så vet vi<br />

att φ kan skrivas som en Fourierserie:<br />

φ(x + iy) = ∑ c(n)<br />

n=0<br />

√ yKiR(2π|n|y)e 2πinx ,<br />

där KiR(y) är en K-Besselfunktion, och om vi approximerar funktionen φ med<br />

en trunkerad Fourierserie kan denna inverteras som en diskret Fouriertransform.<br />

Genom att använda (2) kan vi åstadkomma ett icke-trivialt, välkonditionerat<br />

linjärt ekvationssystem vilket vi kan lösa för att få fram Fourierkoefficienterna<br />

c(n).<br />

Om vi gör detta för ett godtyckligt värde på λ kommer den resulterade<br />

funktionen φ naturligtvis inte vara invariant under Γ. Men om λ är nära ett<br />

riktigt egenvärde så kommer de uträknade Fourierkoefficienterna att vara nära<br />

koefficienterna till en riktig egenfunktion. Med hjälp av vissa egenskaper som<br />

vi vet att en genuin egenfunktion har så bildar vi en “funktional” h(R) som har<br />

egenskapen att den är väldigt nära noll ifall R är nära ett korrekt egenvärde.<br />

Denna funktion används sedan för att lokalisera egenvärden till önskad noggrannhet.<br />

När vi väl har ett egenvärde och en tillhörande uppsättning Fourierkoefficienter<br />

kan vi med hjälp av en ytterligare algoritm (Fas 2) beräkna<br />

ytterligare Fourierkoefficienter på ett väldigt enkelt sätt.<br />

I kapitel 2 studerar vi ett liknande problem, där vi istället för (1) har ekvationen<br />

y 2 φ ′′<br />

xx(z) + φ ′′<br />

yy(z) − iymφ ′ x + λφ(z) = 0,z = x + iy ∈ H , (1’)<br />

där vikten, m, är ett reellt tal, och där villkoret (2) ersätts med<br />

φ(T z) = v(T ) jT (z;m)φ(z), (2’)<br />

där v : Γ → C är ett så kallat multiplikatorsystem, och<br />

jT (z;m) = exp(2πimArg(cz + d)),<br />

för T = <br />

a b<br />

c d och m ∈ R.<br />

Fysikaliskt kan ekvationen (1’) tolkas som att vi nu har en laddad (kvantmekanisk)<br />

partikel som rör sig i ett konstant magnetiskt fält på ytan Γ\H. Den<br />

magnetiska fältstyrkan är proportionell mot m, och multiplikatorsystemet kan<br />

användas för att beskriva flödet genom olika delar av ytan (“hornen”, “handtagen”<br />

etc.) Se t.ex. [2, 1, 11]. (För den som är insatt i fysiken så svarar (1’)<br />

mot vi väljer Landaugaugen, A = c m<br />

y dx för någon lämplig konstant.)<br />

Det visar sig att den metod som infördes i kapitel 1 fortfarande fungerar,<br />

med vissa smärre förändringar. Det (ur beräkningssynpunkt) svåraste problemet<br />

är att Fourierutvecklingen innehåller Whittaker-W- istället för K-Bessel<br />

funktioner. W-funktionen är svårare att beräkna, och vi var tvungna att utveckla<br />

164


en ny algoritm för denna. Algoritmen presenteras i kapitel 4.<br />

I kapitel 3 undersöker vi samma problem som i kapitel 1, men vi betraktar<br />

här även undergrupper Γ ⊆ PSL(2,Z) som inte är kongruensundergrupper.<br />

Den främsta skillnaden mellan kongruens- och icke-kongruensgrupper är<br />

att man inte har lika många “symmetrier” tillgängliga i icke-kongruensfallet.<br />

Detta innebär t.ex. att man inte har så mycket a priori kännedom om vare<br />

sig spektrat eller egenfunktionerna. Man kan därför ställa sig flera frågor<br />

som är mer eller mindre besvarade i kongruensfallet: finns det några “nya”<br />

<strong>Maass</strong>vågfunktioner? hur många finns det i så fall? hur fördelar sig Fourierkoefficienterna?<br />

Uppvisar de några symmmetrier?<br />

Vi bevisar existensen av nya vågfunktioner för grupper Γ som har vissa<br />

specifika egenskaper, vilka på ytorna Γ\H motsvarar olika former av spegelsymmetrier.<br />

(För att systematiskt verifiera dessa egenskaper måste man dock<br />

använda datorer i alla utom de enklaste fallen).<br />

165


Acknowledgements<br />

I would like to express my deep gratitude towards my supervisor, Pr<strong>of</strong>. Dennis<br />

A. Hejhal. He has been my constant guide in this exciting world <strong>of</strong> hyperbolic<br />

geometry, quantum chaos and analytic number theory. I wouldn’t have missed<br />

this for anything! I am also very thankful for all the work he has put in by<br />

pro<strong>of</strong>reading all my manuscripts for this thesis.<br />

I would also like to thank my assistant advisor, Dr. Andreas Strömbergsson<br />

for the many discussions, which have been truly invaluable. My fellow graduate<br />

student, Helen Avelin has also contributed to many <strong>of</strong> these discussions,<br />

and I would like to thank her too, especially in connection with the work in<br />

Chapter 4.<br />

The computational parts <strong>of</strong> my work would not have been possible without<br />

the care <strong>of</strong> the computer system that has been provided by all competent<br />

and friendly computer administrators here at the department and at NSC in<br />

Linköping.<br />

When widening the area <strong>of</strong> gratefulness, I would also like to thank all my<br />

co-workers at the department for being around and providing a nice atmosphere<br />

for me to work and drink c<strong>of</strong>fee in. I mention no one in particular,<br />

since that’s the only way I won’t forget anyone.<br />

Close by the department, I should also thank Möbius (the student organization)<br />

for providing me with a constant supply <strong>of</strong> cc-light for the last 12 years<br />

or so.<br />

Thanks to all my friends and acquaintances in Uppsala for providing nonmathematical<br />

reasons to stay here.<br />

While I am at it, I would also like to thank my parents and family for this<br />

(once in a lifetime) opportunity to exist.<br />

And, finally, I would like to thank ............... (insert your name here if you<br />

feel left out).<br />

167


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