Master Dissertation
Master Dissertation
Master Dissertation
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<strong>Master</strong> <strong>Dissertation</strong><br />
The Method of Epstein and Glaser<br />
Author: Asger Jacobsen Supervisor: J.P. Solovej<br />
21st August 2005 Temp. version
Contents<br />
1 Introductory Theory 6<br />
1.1 Notation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 6<br />
1.2 Formal Power Series . . . . . . . . . . . . . . . . . . . . . . . 8<br />
1.3 Affine Transformations of Distributions . . . . . . . . . . . . 12<br />
2 The Poincaré Group 14<br />
2.1 The Lorentz Group . . . . . . . . . . . . . . . . . . . . . . . . 14<br />
2.2 The Poincaré Group . . . . . . . . . . . . . . . . . . . . . . . 16<br />
2.3 Spinor Representations of the Lorentz Group . . . . . . . . . 17<br />
3 The Scattering Matrix 19<br />
4 The Mathematical Setting of QFT 22<br />
4.1 The Fock Space . . . . . . . . . . . . . . . . . . . . . . . . . . 22<br />
4.2 The Wightman Axioms . . . . . . . . . . . . . . . . . . . . . 25<br />
5 The Method of Epstein and Glaser 27<br />
5.1 Introduction . . . . . . . . . . . . . . . . . . . . . . . . . . . . 27<br />
5.2 Example - In the Hilbert Space Setting of Quantum Mechanics 36<br />
5.3 Splitting of Distributions . . . . . . . . . . . . . . . . . . . . . 38<br />
6 Regularly Varying Functions 41<br />
7 Splitting of Numerical Distributions 49<br />
7.1 The Singular Order of a Distribution . . . . . . . . . . . . . . 49<br />
7.2 Case I: Negative Singular Order . . . . . . . . . . . . . . . . . 54<br />
7.2.1 Existence . . . . . . . . . . . . . . . . . . . . . . . . . 54<br />
7.2.2 Uniqueness . . . . . . . . . . . . . . . . . . . . . . . . 59<br />
7.3 Case II: Positive Singular Order . . . . . . . . . . . . . . . . . 64<br />
8 Application to QED 66<br />
8.1 Using the Game Plan . . . . . . . . . . . . . . . . . . . . . . 66<br />
8.2 The Adiabatic Limit . . . . . . . . . . . . . . . . . . . . . . . 69<br />
ii
9 The Microlocal Approach - A Condition on the Wave Front<br />
Set. 71<br />
iii
Abstract<br />
The main purpose of this thesis is to prove that the causality axiom allows<br />
a well-defined perturbation theory of quantum electrodynamics. This<br />
will be done by using the method of Epstein and Glaser. The thesis introduces<br />
the theory of formal power series and regularly varying functions, both<br />
tools needed in the main proofs. Basic theory of the scattering matrix, the<br />
Poincaré Group and its representations will be introduced. We will use the<br />
concept of singular order to investigate the splitting of distributions into an<br />
advanced and a retarded part and show when they exist and when they are<br />
unique. Applications and the adiabatic limit will be discussed. The thesis<br />
is concluded with a consideration of the microlocal approach to the method<br />
of Epstein and Glaser and by showing that the translations invariance can<br />
be substituted by a condition on the wave front set.<br />
Sammenføjning<br />
Hovedform˚alet med dette speciale er at vise at man ved hjælp af kausalitetsaxiomet<br />
kan indføre en veldefineret pertubationsteori for kvanteelektrodynamik.<br />
Dette gøres ved hjælp af Epstein og Glasers metode. Specialet introducerer<br />
teorierne for formelle potensrækker og regulært varierende funktioner,<br />
som begge skal bruges under hovedbeviserne. Basal teori for spredningsmatrice,<br />
Poincaré gruppen og dens repræsentationer vil blive introducerede.<br />
Vi vil bruge singulær orden til at undersøge opdelingen af distributioner<br />
i en fremtids og fortids del og vise, hvorn˚ar de eksisterer, og hvorn˚ar<br />
de er entydige. Anvendelser og den adiabatiske grænse vil blive diskuteret.<br />
Specialet afsluttes med at overveje den mikrolokale tilgangsvinkel til Epstein<br />
og Glasers metode og med at vise at translations varians kan afløses af en<br />
betingelse p˚a bølgefront mængden. Specialet er skrevet p˚a engelsk.<br />
1
Preface<br />
Overview<br />
This dissertation is about the causal approach to finite quantum electrodynamics<br />
(QED). The word causal refers to the main assumption on which the<br />
method we will introduce is based. In short it is the statement that what<br />
happens at time t > s doesn’t influence what happens at earlier times t < s.<br />
The word finite means that there do not arise any ultraviolet divergences.<br />
The traditional formalism of quantum field theory (QFT) is known to<br />
have problems with ultraviolet divergences. These occur because of misuse<br />
of the mathematics. The problem lies in the fact that QFT deals with distributions<br />
instead of functions and that the rules of calculus for these are<br />
not the same. One cannot simply multiply a distribution by a discontinuous<br />
step-function and other simple rules of calculus like multiplication of distributions<br />
can be a very complicated matter. In fact, this is the main problem<br />
of [4].<br />
In practice the ultraviolet divergencies are dealt with by regularization<br />
and renormalization FIXME: explain these. But instead of using these firstaid<br />
features on an ill-defined theory one could introduce QFT using the<br />
mathematics in the right way from the beginning. This is the basic aim of<br />
this dissertation.<br />
Structure<br />
Chapter 1<br />
In Chapter 1 I introduce the tools needed in our later work. Section 1.2<br />
is about formal power series which is an abstraction of the usual power<br />
series used in calculus. A generalization which lets us use most of the wellknown<br />
machinery from calculus to settings which have no natural notion<br />
of convergence. This will become useful later on when we introduce the<br />
scattering matrix as a formal power series.<br />
In section 1.3 we will introduce Affine Transformations of distributions.<br />
2
Chapter 2<br />
The Poincaré Group will be introduced in Sections 2.1 and 2.2 in Chapter 2.<br />
Further, in Section 2.3 the spinor representation of the Lorentz Group will<br />
be presented.<br />
Chapter 3<br />
In Chapter 3 we will take a look at the key player of this dissertation: The<br />
Scattering Matrix. In this chapter we introduce it in the well-known setting<br />
of quantum mechanics.<br />
Chapter 4<br />
In Chapter 4 we consider the setting of QFT. In Section 4.1 we introduce<br />
the Fock space and in Section 4.2 we take a look at the Wightman axioms.<br />
Chapter 5<br />
In Chapter 5 we will finally introduce the method of Epstein and Glaser.<br />
Section 5.1 is about the main ideas and a game plan is sketched. In Section<br />
5.2 the method is applied to the of quantum mechanics which is based<br />
on the well-known theory of Hilbert spaces. In Section 5.3 we will look at<br />
the splitting of distributions into an advanced and a retarded part suggested<br />
in the game plan. We focus on the support properties.<br />
Chapter 6<br />
In Chapter 6 we will leave the theory for a while introducing an important<br />
tool for the splitting of numerical distributions in Chapter 7.<br />
Chapter 7<br />
In Section 7.1 of this chapter we develop some tools for studying the splitting<br />
of numerical distributions. The notion of singular order of a distribution will<br />
be a tool to categorize the type of splitting. Thus we will in Sections 7.2<br />
and 7.3 look at the splitting of of distributions of positive and negative<br />
singular order, separately. We will find out that the splitting may be done<br />
by multiplication of a step-function (this will be done in a technical manner)<br />
in the case of negative order. Further in this case the splitting is unique,<br />
which will be the topic of Subsection 7.2.2. The splitting in the case of<br />
positive singular order turns out to be a more complicated matter and only<br />
works for a special class of distributions. Further the splitting is not unique.<br />
3
Chapter 8<br />
In Chapter 8 we will turn our attention to applications. In Section 8.1 we<br />
find an expression for the difference distribution which is essential in the<br />
splitting process. The expression can be considered a correspondent to the<br />
Feynman Rules. We find the the first term describes photon exchange.<br />
In Section 8.2 we will give an idea of how to take the adiabatic limit.<br />
Chapter 9<br />
In this chapter we consider the microlocal approach to the method of Epstein<br />
and Glaser. The idea is to generalize the theory to manifolds. As<br />
translations are substituted by parallel transport on manifolds the property<br />
of translational invariance has to be substituted by a condition on the<br />
smoothness. We find a condition on the wave front set which makes the<br />
terms of the scattering matrix well-defined.<br />
Thanks<br />
I would like to thank my supervisor Jan Philip Solovej for his help and<br />
support during the writing of this project. Further I would like the numerous<br />
people who’s work and research this dissertation is based on; especially<br />
Scharf, Fredenhagen, and of course Epstein and Glaser. Finally, I will also<br />
thank my family who has had to bear over with my during the final work<br />
under time pressure.<br />
4
Prerequisites<br />
Fixme: Reevaluate the needs as the theory develops.<br />
This is a dissertation for the masters degree in mathematics at the University<br />
of Copenhagen. It is thus the conclusion of 5 years of study. The<br />
thesis is meant to rest on the mathematical fundament the student has obtained<br />
during those years. Clearly, I will not use all the mathematical tools<br />
I am equipped with from the wide variety of disciplines, I been acquainted<br />
with through my studies. It would be more correct to say that it is a narrow<br />
specialization of some of the knowledge I have gained. It is thus difficult<br />
to point out what the prerequisites are, exactly. I will never the less try to<br />
narrow down the fundament on which this thesis rests.<br />
First of all I will emphasize that this thesis is concerned with quantum<br />
electrodynamics. That is, it is concerned with problems arising from both<br />
quantum mechanics and general relativity. Although I will only be concerned<br />
with a few principles of the two, the whole idea and machinery of them are<br />
luring in the background. I have through my studies read several books on<br />
the subjects most important probably being the basic theory to be found in<br />
[7], [8] and [1]. Note that Stone’s Theorem will be taken for granted and <br />
will be normalized.<br />
Also a good deal of functional analysis will be taken for granted. At least<br />
the content of [5]. Especially the theory of distributions will be imperative<br />
the thesis. I will, in some sense, consider the project [4], which I wrote with<br />
my fellow student Morten Bakkedal, as a part of this thesis. The project<br />
covers a good deal of the theory of distributions and wavefront sets resting<br />
on the lecture notes on distributions by Gerd Grubb [2].<br />
5
Chapter 1<br />
Introductory Theory<br />
1.1 Notation<br />
I will mostly adopt the notation used in [4], but for convenience there will be<br />
some differences in the notation of a distribution. Whereas we in [4] used to<br />
write a distribution u acting on a test-function φ in the functional manner<br />
u(φ) I will in this project use the “inner product” notation 〈u, φ〉.<br />
Also note that I will use the definition of the fourier transform found in<br />
[4] rather than the one found in [6]. This will only result in some differences<br />
in the factors.<br />
I will occasionally use Einstein notation along with the convention that<br />
indices from the Greek alphabet take the values 0, 1, 2, 3 and indices from<br />
the Latin alphabet take the values 1, 2, 3.<br />
By a t-dependent operator O(t) : E → F we mean an indexed family of<br />
operators {O(t)}t∈R all from E to F .<br />
Given a t-dependent operator O(t) : E → F . We write O = s- limt→∞ O(t)<br />
if it converges in the strong operator topology to an operator O. That is, if<br />
Oφ = lim<br />
t→∞ (O(t)φ) for all φ ∈ E.<br />
Beware of the many meanings of the symbol δ. Usually it means the<br />
Dirac delta-distribution, but it is also used as a small real, the Kronecker<br />
symbol δi,j and the diagonal map. It should be transparent from the context<br />
which meaning the symbol has.<br />
For convenience I have listed some general notation on the next page:<br />
6
General notation.<br />
Let α = (α1, . . . , αn) ∈ N n 0 be a multi-index and x = (x1, . . . , xn) ∈ R n a<br />
point, then<br />
|α| = α1 + α2 + . . . + αn.<br />
α! = (α1, α2, . . . , αn)! = α1!α2! · · · αn!.<br />
x α = x α1<br />
1<br />
· · · xαn<br />
n<br />
∂j = ∂/∂xj for j = 1, . . . , n.<br />
∂ α = ∂ α1<br />
1<br />
· · · ∂αn<br />
n .<br />
Dj = −i∂j for j = 1, . . . n.<br />
f (α) (x) = ∂ α f(x).<br />
7
1.2 Formal Power Series<br />
The concept of formal power series will be of much importance to this<br />
project. They make it possible to employ much of the analytical machinery<br />
of usual power series but in settings which have no natural notion of<br />
convergence.<br />
Let R be a commutative ring with identity. We call the elements of R the<br />
scalars. Let R N be the set of all infinite sequences in R. Now we define<br />
addition and multiplication of two such sequences in the usual way, that is,<br />
and the Cauchy product<br />
(an) + (bn) = (an + bn),<br />
n<br />
(an) × (bn) =<br />
k=0<br />
akbn−k<br />
<br />
= <br />
h+k=n<br />
ahbk,<br />
which in a sense is discrete convolution. Note that addition and<br />
multiplication are well-defined operations because each term only depends<br />
on a finite number of terms from the original sequences.<br />
Theorem 1.1. (R N , +, ×) is a commutative ring with the multiplicative<br />
identity 1 := (1, 0, 0, . . .) and additive identity 0 := (0, 0, . . .).<br />
Proof. The only non-trivial property is multiplicative associativity. We<br />
need to prove that<br />
((a × b) × cn) = (a × (b × c)n).<br />
This follows from<br />
<br />
((a × b) × cn) =<br />
<br />
(a × b)hck =<br />
<br />
and<br />
h+k=n<br />
<br />
(a × (b × c)n) =<br />
r+h=n<br />
r+s+k=n<br />
<br />
ar(b × c)h =<br />
r+s+k=n<br />
arbsck<br />
arbsck<br />
It is easy to see from the definition of the Cauchy product that if we define<br />
X = (0, 1, 0, 0, . . .) then every element of R N on the form<br />
(a0, . . . , aN, 0, . . .) can be written on the form<br />
N<br />
anX n .<br />
n≥0<br />
8<br />
<br />
<br />
.
Theorem 1.2. Let d : R N × R N → R be defined by<br />
d((an), (bn)) = 2 −k , where k ∈ N0<br />
is the smallest index such that ak = bk. If no such k exists we define<br />
d((an), (bn)) = 0. Then (R N , d) is a metric space.<br />
Proof. The only non-trivial part is the triangle inequality.<br />
Let d(a, b) = 2−k1 −k2 , d(b, c) = 2<br />
exist.<br />
−k3 and d(a, c) = 2 assuming k1, k2 and k3<br />
Say, k3 ≥ k1, then 2−k3 −k1 −k1 −k2 ≤ 2 ≤ 2 + 2 , hence we may assume<br />
= bk3 and bk3 = ck3 thus ak3 = ck3 ,<br />
k3 < ki for i = 1, 2. Then ak3<br />
contradicting the definition of k3.<br />
If k3 doesn’t exist the inequality is obvious. Say, k1 doesn’t exist, then<br />
(an) = (bn) hence k2 = k3 if they exist. If not, then (an) = (cn).<br />
Theorem 1.3. The operations of addition and multiplication are<br />
continuous.<br />
Proof. Say, fn → f and gn → g then we need to prove that<br />
fn + gn → f + g, that is d(fn + gn, f + g) → 0 for n → ∞.<br />
We need to prove that given ɛ > 0 there exists a δ such that<br />
max{d(fn, f), d(gn, g)} < δ ⇒ d(fn + gn, f + g) < ɛ.<br />
Given ɛ > 0 find k such that 2 −k < ɛ. Let δ = 2 −k /2 = 2 −(k+1) , then<br />
|(f − fn)h| + |(g − gn)h| = 0, for all h ≤ k + 1. (1.1)<br />
Hence since 2 −(k+1) + 2 −(k+1) = 2 −k<br />
d(fn + gn, f + g) < 2 −k < ɛ.<br />
Note that equation (1.1) also implies |(fngn)h| − |(fg)h|=0 for h ≤ k + 1,<br />
hence multiplication is also continuous.<br />
The cases where no k exists are obvious.<br />
This leads to the following result.<br />
Corollary 1.4. (RN , d) is a topological ring and<br />
(an) = <br />
anX n .<br />
n≥0<br />
Note that in this ring convergence is absolute, in fact any rearrangements<br />
of the series converges to the same limit.<br />
Definition 1.5. We denote the topological ring (R N , d) by R[[X]] and call<br />
it the ring of formal power series.<br />
9
Theorem 1.6. Let anX n be a formal power series. Then anX n has<br />
a unique inverse if and only if a0 has a multiplicative inverse in R.<br />
Proof. Say, bnX n is the inverse of anX n . Then by the definition of<br />
multiplication we must have<br />
(i) a0b0 = 1<br />
(ii) a0b1 + a1b0 = 0<br />
(iii) a0b2 + a1b1 + a2b0 = 0.<br />
and so on. . .<br />
Now it is clear from (i) that if the inverse exists then a0 must also be<br />
invertible.<br />
If on the other hand a0 ∈ R ∗ , the invertible elements of R, then the inverse<br />
of anX n is given by bnX n , where b0 = a −1<br />
0 and<br />
bn = a −1<br />
0 (−a1bn−1 − . . . − anb0), for n ≥ 1.<br />
Say, both bnX n and cnX n are inverse to anX n . Then a0b0 = a0c0,<br />
that is b0 = c0, since a0 = 0. Assume bk = ck, for all k < n, then<br />
a0bn + a1bn−1 + · · · + anb0 = a0cn + a1cn−1 + · · · + anc0<br />
By the induction assumption bk = ck for all k < n hence a0bn = a0cn and<br />
thus bn = cn.<br />
The geometric series formula is valid in R[[X]].<br />
Theorem 1.7. Let anX n = 1 be a formal power series, then<br />
<br />
anX n m<br />
= (1 − anX n ) −1 .<br />
Proof. Let g = <br />
anX n<br />
m . Then<br />
<br />
1 − anx n<br />
<br />
g = 1 − anx n<br />
as wanted.<br />
lim<br />
M→∞<br />
m=0<br />
M <br />
anx n ) m<br />
<br />
= lim 1 −<br />
M→∞<br />
anx n M<br />
m=0<br />
<br />
= 1 − anx n<br />
<br />
+ anx n −<br />
= 1.<br />
anx n ) m<br />
Setting an = δ1,n we get the usual geometric series formula.<br />
10<br />
anx n 2<br />
+ · · ·
Corollary 1.8.<br />
∞<br />
n=0<br />
X n = (1 − X) −1 .<br />
Note the following nice properties of R[[X]].<br />
Lemma 1.9. The topology on R[X] is equal to the product topology on<br />
R N , where R is equipped with the discrete topology.<br />
Proof. Let B(a, ɛ) be an open ball in R[[X]], then<br />
B(a, ɛ) = {b ∈ R[[X]]|d(a, b) < ɛ}<br />
= {b ∈ R[[X]]|d(a, b) < 2 −k , where 2 −k ≤ ɛ < 2 −k+1 }<br />
= {b ∈ R[[X]]|ah = bh for all 0 ≤ h < k}<br />
=<br />
k−1 <br />
h=0<br />
π −1<br />
h ({ah}),<br />
where {ah} is open in the discrete topology.<br />
If on the other hand<br />
B = π −1<br />
β1<br />
(Uβ1 ) ∩ π−1(Uβ2<br />
) ∩ · · · ∩ π−1(Uβn<br />
),<br />
β2<br />
is a basis element of R N . Let a ∈ B. Define m = max{β1, . . . , βn}. Then<br />
a ∈ B a, 2 −(m+1) which is clearly a subset of B.<br />
Proposition 1.10. The metric space (R[[X]], d) is<br />
1. Complete.<br />
2. Compact if and only if R is finite.<br />
Proof. Completeness: Let (an) be a cauchy sequence in R[[X]]. Then we<br />
know<br />
∀k∃N ∈ N∀n, m ≥ N : d(an, am) < 2 −k .<br />
Hence for each h, πh(an) converges to, say, bh. Then (an) converges to the<br />
formal power series b = <br />
h bhX h .<br />
Compactness: If R is finite R[[X]] is compact by Tychonoff’s theorem.<br />
The topology of R is the discrete topology, hence if it is infinite, <br />
a∈R a is<br />
an open cover that doesn’t contain a finite subcollection covering R.<br />
11<br />
βn
1.3 Affine Transformations of Distributions<br />
Let A = (aij) be an invertible n × n matrix and f a continuous function on<br />
R n . We define a linear continuous composition ⋆ by<br />
A ⋆ f(x) = f(Ax), x ∈ R n .<br />
Let’s look at the distribution induced by f acting on a test-function φ.<br />
Letting T (x) = A −1 x, where A −1 = (ãij), by change of variables theorem<br />
<br />
〈A ⋆ f, φ〉 =<br />
<br />
=<br />
<br />
=<br />
R n<br />
since (A −1 x)i = n<br />
j=1 ãijxj and<br />
T ′ (x) =<br />
⎛<br />
⎜<br />
⎝<br />
∂(A −1 x)1<br />
∂x1<br />
.<br />
∂(A −1 x)n<br />
∂x1<br />
f(Ax)φ(x)dx<br />
f(T (Ax))φ(T (x))| det T ′ (x)|dx<br />
Rn f(x)φ(A −1 x)| det A −1 |dx<br />
= | det A| −1 〈f, A −1 ⋆ φ〉,<br />
. . .<br />
. ..<br />
. . .<br />
∂(A −1 x)1<br />
∂xn<br />
.<br />
∂(A −1 x)n<br />
∂xn<br />
⎞<br />
⎟<br />
⎠ =<br />
⎛<br />
⎞<br />
ã11 . . . ã1n<br />
⎜<br />
⎝<br />
.<br />
. ..<br />
⎟<br />
. ⎠ = A −1 .<br />
ãn1 . . . ãnn<br />
Definition 1.11. Let u ∈ D ′ (R n ) and let A be a real invertible n × n<br />
matrix. Then the distribution A ⋆ u is defined by<br />
for φ ∈ C ∞ 0 (Rn ).<br />
〈A ⋆ u, φ〉 = | det A| −1 〈u, A −1 ⋆ φ〉, (1.2)<br />
We should check that this actually defines a distribution.<br />
Linearity:<br />
〈A ⋆ u, rφ + sψ〉 = | det A| −1 〈u, A −1 ⋆ (rφ + sψ)〉<br />
Continuity: Assume φn → φ, then<br />
= | det A| −1 〈u, rA −1 ⋆ φ + sA −1 ⋆ ψ〉<br />
= r| det A| −1 〈u, A −1 ⋆ φ〉 + s| det A| −1 〈u, A −1 ⋆ φ〉<br />
〈A ⋆ u, φn〉 = | det A| −1 〈u, A −1 ⋆ φn〉 → | det A| −1 〈u, A −1 ⋆ φ〉,<br />
since u and ⋆ are continuous.<br />
12
Proposition 1.12. We state some often occurring affine transformations.<br />
In the following u ∈ D ′ (R n ) and φ ∈ C(R n ).<br />
1. Reflection: 〈u(−x), φ(x)〉 = 〈u(x), φ(−x)〉.<br />
2. Scaling: 〈u(tx), φ(x)〉 = t −n 〈u(x), φ(x/t)〉 = t −n 〈u(x), λtφ(x)〉, letting<br />
λtφ(x) := φ(x/t).<br />
Proof. 1. and 2. follow directly from (1.2) by choosing A = −I and A = tI<br />
respectively.<br />
Note that λt : S(R n ) → S(R n ) is continuous. This follows since<br />
λtφα,β = sup |x α D β (λt(φ(x)))| = sup |x α t −|β| (D β φ)(x/t)|<br />
Further we define<br />
= t |α|−|β| sup |(x/t) α (D β φ)(x/t)|<br />
= t |α|−|β| φα,β.<br />
Definition 1.13. Let u ∈ D ′ (R n ), h ∈ R n and φ ∈ C(R n ). We define the<br />
translation τh of u by<br />
〈u(x − h), φ(x)〉 = 〈τhu, φ〉 = 〈u, τ−hφ〉 = 〈u(x), φ(x + h)〉.<br />
Clearly this defines a distribution.<br />
Recall that a function f is said to be homogeneous of order λ ∈ C if<br />
f(tx) = t λ f(x) for all t > 0 and x ∈ R n . We need to define the similar<br />
property for distributions.<br />
Definition 1.14. A distribution u is said to be homogeneous of degree<br />
λ ∈ C if<br />
〈u(tx), φ(x)〉 = 〈t λ u(x), φ(x)〉,<br />
for all t > 0 and φ ∈ C ∞ 0 .<br />
13
Chapter 2<br />
The Poincaré Group<br />
2.1 The Lorentz Group<br />
We define the Lorentz Metric as the bilinear form on R4 seen as a vector<br />
space:<br />
〈x, y〉 = x 0 y 0 − x 1 y 1 − x 2 y 2 − x 3 y 3<br />
for all x, y ∈ R 4 . (2.1)<br />
Obviously it is symmetric and non-degenerate in the sense that there for<br />
all x = 0 exists a y such that 〈x, y〉 = 0. But it is not positive definite.<br />
Letting g be the Metric Tensor i.e.<br />
<br />
−1 0T g = (gµν) =<br />
, where 0<br />
0 I3<br />
T = (0, 0, 0).<br />
Equation (2.1) becomes<br />
〈x, y〉 = x T gy = gµνx µ y ν<br />
The Minkowski Space M is the vector space R 4 endowed with the Lorentz<br />
metric.<br />
While y µ denotes the components of the vector y in M, yµ denotes the<br />
components of the corresponding linear form y ′ given by y ′ (x) = 〈y, x〉, for<br />
all x. In the dual space with the canonical basis the components of the<br />
linear form y ′ are (y 0 , −y 1 , −y 2 , −y 3 ).<br />
Definition 2.1. A Lorentz transformation of R 4 is a linear map<br />
Λ : R 4 → R 4 satisfying<br />
〈Λx, Λy〉 = 〈x, y〉 for all x, y ∈ R 4 . (2.2)<br />
Since e T i Aej = Aij for a matrix A it is clear that (2.2) is equivalent to<br />
Λ T gΛ = g. (2.3)<br />
14
Clearly the composition of Lorentz transformations is itself a Lorentz<br />
transformation. And since<br />
1 = − det g = − det(Λ T gΛ) = − det Λ det g det Λ = (det Λ) 2 ,<br />
we must have det Λ = ±1. Thus all Lorentz transformations are invertible<br />
and the set of all Lorentz transformations form group called the Lorentz<br />
Group L, also denoted O(3, 1).<br />
Further we denote the group of Lorentz transformations Λ with det Λ = 1<br />
by L+ = SO(1, 3). From equation (2.3) we get 10 independent quadratic<br />
equations for the components of Λ. The first is<br />
Which implies that either<br />
When Λ 0 0<br />
(Λ 0 0) 2 − (Λ 1 0) 2 − (Λ 2 0) 2 − (Λ 3 0) 2 = 1. (2.4)<br />
Λ 0 0 ≥ 1 or Λ 0 0 ≤ −1. (2.5)<br />
≥ 1 time is not reversed and we call the subgroup defined by<br />
the Proper Lorentz Group.<br />
L ↑<br />
+ := {Λ ∈ L+|Λ 0 0 ≥ 1},<br />
The Lorentz group is normally split into 4 disjoint classes. Namely, L ↑<br />
+<br />
and the 3 following<br />
Name det Λ Λ 0 0<br />
L ↑<br />
− −1 ≥ +1<br />
L ↓<br />
− −1 ≤ −1<br />
L ↓<br />
+ +1 ≤ −1<br />
Further we mention 3 important discrete Lorentz transformations, namely:<br />
I, the identity, P = g, space inversion or the parity transform, and<br />
T = −g, time inversion.<br />
There are 3 types of vectors in M. A vector x is said to be time-like if<br />
x 2 > 0, space-like if x 2 < 0 and light-like if x 2 = 0. Since x 2 is constant<br />
under Lorentz transformation the categories stay the same under the<br />
transformation.<br />
Information cannot propagate faster than the speed of light. Equivalently,<br />
two events at x µ and y ν cannot influence each other if they are separated<br />
by a space-like distance,<br />
(x − y) 2 < 0.<br />
They have no causal relation. Causality plays a key role in this thesis.<br />
15
Definition 2.2. Given some point x the forward cone of x is defined by<br />
and the backwards cone of x is<br />
V + (x) = {y|(y − x) 2 ≥ 0 and y 0 ≥ x 0 },<br />
V − (x) = {y|(y − x) 2 ≥ 0 and y 0 ≤ x 0 }.<br />
The n-dimensional generalizations are<br />
Γ ± n (x) = {(x1, . . . , xn)|xj ∈ V ± (x) and ∀j = 1 = 1, . . . , n}.<br />
That (y − x) 2 ≥ 0 just means that they have to be causally connected.<br />
Later we will need a symbol to distinguish points in time.<br />
Definition 2.3. Let A, B ⊂ M. If for all x ∈ A and y ∈ B we have<br />
x 0 < y 0 , we write A < B.<br />
2.2 The Poincaré Group<br />
Definition 2.4. Let Λ ∈ L and a ∈ R 4 . By a Poincaré transformation<br />
Π : R 4 → R 4 we mean Π(x) = Λx + a, and we write Π = (a, Λ).<br />
The set P of all Poincaré transformations form a group under the<br />
composition law<br />
(a1, Λ1)(a2, Λ2) = (a1 + Λ1a2, Λ1Λ2).<br />
We see that the Poincaré Group P is the semidirect product of L and the<br />
group of space-time translations (R 4 , +), that is<br />
P = R 4 ⊙ L<br />
Further we define the Proper Poincaré Group as<br />
P ↑<br />
+ = R4 ⊙ L ↑<br />
+<br />
16
2.3 Spinor Representations of the Lorentz Group<br />
We want to extend the spinor representation to four dimensions. 1 Doing<br />
this we add the unit matrix<br />
<br />
1 0<br />
σ0 = , (2.6)<br />
0 1<br />
to the usual Pauli matrices:<br />
σ1 =<br />
0 1<br />
1 0<br />
<br />
, σ2 =<br />
0 −i<br />
i 0<br />
<br />
1 0<br />
, σ3 =<br />
0 −1<br />
<br />
. (2.7)<br />
These form a basis of H(2), the vector space over R of all complex<br />
Hermitian 2 × 2 matrices. Actually, given a 2 × 2 matrix H ∈ H(2) we get<br />
by direct calculation that H can be represented by a vector x ∈ R4 in the<br />
following way:<br />
H := σ(x) = x µ σµ = 1<br />
3<br />
tr(Hσµ)σµ<br />
(2.8)<br />
2<br />
The map x ↦→ σ(x) defined in this way is indeed an isomorphism of R 4<br />
onto H(2).<br />
Fixme: should this be written out?<br />
Further we can define another isomorphism by<br />
µ=0<br />
x ↦→ σ ′ (x) = x 0 σ0 − x k σk =<br />
3<br />
xµσµ.<br />
Now consider the transform of σ(x) with a complex 2 × 2 matrix A:<br />
σ(x) ↦→ Aσ(y)A ∗ .<br />
Clearly σ(y) := Aσ(x)A ∗ ∈ H(2). And from equation (2.8) we see that its<br />
components in the basis of Pauli Matrices are given by<br />
y ν = 1<br />
2 tr(σ(y)σν) = 1<br />
2 tr(Aσ(x)A∗σν) = 1<br />
2<br />
µ=0<br />
3<br />
µ=0<br />
tr(AσµA ∗ σν)x µ .<br />
In this way we get a linear map ΛA : x → y of R 4 into it self. Note that<br />
Hence if det A = 1, then<br />
σ(x) = x µ <br />
x0 + x3 x1 − ix2 σµ =<br />
x 1 + ix 2 x 0 − x 3<br />
<br />
.<br />
1 The two-component spinor formalism is treated in [8] chapter 3.2-3.3.<br />
17
〈y, y〉 = det σ(y) = det A det σ(x) det A ∗ = det σ(x) = 〈x, x〉.<br />
Clearly the parallelogram identity holds for the Lorentz metric 2 , thus<br />
〈x, y〉 =<br />
= det σ(x) + det σ(y) − det σ(x) − det σ(y)<br />
〈x + y, x + y〉 − 〈x, x〉 − 〈y, y〉<br />
2<br />
.<br />
2<br />
From this we see that 〈ΛAx, ΛAy〉 = 〈x, y〉, thus it is a Lorentz<br />
transformation.<br />
Since<br />
fixme: make sure the following is correct<br />
Hence<br />
ΛB(x) = y ⇔ σ(y) = Bσ(x)B ∗ and ΛA(y) = z ⇔ σ(z) = Aσ(y)A ∗ ,<br />
ΛAB(x) = z ⇔ ABσ(x)B ∗ A ∗ = Aσ(y) = σ(z).<br />
ΛAB = ΛAΛB.<br />
In this way we can construct a representation Λ : A ↦→ ΛA of SL(2, C), the<br />
group of all complex 2 × 2 matrices with determinant 1, onto L ↑<br />
+ . We call<br />
SL(2, C) the spinor representation. The vectors in the representation space<br />
C2 are called spinors.<br />
Note that the kernel of the representation Λ is<br />
ker Λ = Λ −1 ({I4}) = {A ∈ SL(2, C)|AHA ∗ for all H ∈ H(2)}.<br />
Letting H = I2 we see that A ∈ ker Λ must be unitary and A ∈ ker Λ if and<br />
only if AH = HA for all H ∈ H(2). Hence A = −I2 or A = I2 and<br />
L ↑<br />
+ ∼ = SL(2)/{I2, −I2}.<br />
Hence Λ(A) = Λ(−A) and SL(2, C) is a two-valued representation of L ↑<br />
+ .<br />
2 The proof is exactly the same as the usual one for inner product spaces.<br />
18
Chapter 3<br />
The Scattering Matrix<br />
As the notion of the scattering matrix, or the S-matrix for short, plays the<br />
central role in scattering theory I will give a brief introduction to the<br />
development of it in the case of quantum mechanics. Later our main aim<br />
will be to construct the S-matrix of quantum electrodynamics, QED, by<br />
perturbation theory. First a definition:<br />
Definition 3.1. Let H be a Hilbert space. Then a two-parameter family<br />
U(s, t), (s, t) ∈ R 2 of bounded operators on H is called a unitary<br />
propagator, if the following is satisfied:<br />
(1) U(s, t) is unitary for all s, t ∈ R,<br />
(2) U(t, t) = 1 for all t ∈ R,<br />
(3) U(r, s)U(s, t) = U(r, t) for all r, s, t ∈ R,<br />
(4) The map (s, t) ↦→ U(s, t)ψ is continuous for all ψ ∈ H.<br />
Consider a quantum mechanical system described by the Hamiltonian<br />
H = H0 + V (t), where H0 is the free Hamiltonian and V is the<br />
time-dependent interaction. Then the time evolution is given by a unitary<br />
propagator in the sense that<br />
ψ(t) = U(t, s)ψ(s) 1 . (3.1)<br />
Definition 3.2. Let U(s, t) be a unitary propagator then we define the<br />
wave operators as follows<br />
Win = s- lim<br />
t→−∞ U(t, 0)∗ e −iH0t = s- lim<br />
t→−∞ U(0, t)e−iH0t .<br />
Wout = s- lim<br />
t→∞ U(t, 0) ∗ e −iH0t = s- lim<br />
t→∞ U(0, t)e −iH0t .<br />
Provided the strong limits exist.<br />
1 We remember that if H is time independent then then time evolution is given by the<br />
unitary transform ψ(t) = e −iHt ψ(t0).<br />
19
Note that<br />
〈Woutφ, ψ〉 = 〈 lim<br />
t→∞ U(0, t)e −iH0t φ, ψ〉<br />
= lim<br />
t→∞ 〈U(0, t)e −iH0t φ, ψ〉<br />
= lim<br />
t→∞ 〈φ, e iH0t U(0, t) ∗ ψ〉<br />
= 〈φ, lim<br />
t→∞ e iH0t U(t, 0)ψ〉.<br />
Hence W ∗ out = s- limt→∞ e iH0t U(t, 0). Further<br />
W ∗ outWinψ = lim<br />
t→∞ lim<br />
s→−∞ eiH0t U(t, 0)U(0, s)e −iH0s ψ<br />
This leads us to the following definition.<br />
= lim<br />
t→∞ lim<br />
s→−∞ eiH0t U(t, s)e −iH0s ψ.<br />
Definition 3.3. The scattering matrix is defined as follows<br />
S = W ∗ outWin = s- lim<br />
t→∞ s- lim<br />
s→−∞ eiH0t U(t, s)e −iH0s .<br />
The physical meaning of the scattering matrix is now apparent. A<br />
normalized initial asymptotic state ψ considered at time t = 0, say, is first<br />
transformed to s = −∞ by free dynamics, then it is evolved from −∞ to<br />
t = ∞ by full interacting dynamics and finally it is transformed back from<br />
∞ to t = 0 again by free dynamics. Thus Sψ is in fact the outgoing<br />
scattering state transformed to t = 0 by free dynamics. The probability for<br />
a transition form ψ to φ is given by<br />
P (ψ → φ) = |〈φ, Sψ〉| 2 .<br />
Now ψ(t) as given by equation (3.1) is the solution to the Schrödinger<br />
equation<br />
i d<br />
dt ψ(t) = (H0 + V (t))ψ(t).<br />
If we go over to the interaction picture 2 by substituting φ = e iH0t ψ, we get<br />
i d d<br />
φ(t) = i<br />
dt dt (eiH0tψ(t)) = −H0e iH0t iH0t<br />
ψ(t) + e (H0 + V (t))ψ(t)<br />
= e iH0t −iH0t<br />
V e φ(t)<br />
= ˜ V (t)φ(t).<br />
FIXME: Make sure you can differentiate like this.<br />
2 See [8] page 318.<br />
20
Thus we note that the scattering matrix is in fact just the limit of the time<br />
evolution in the interaction picture. That is<br />
S = lim lim Ũ(t, s).<br />
t→∞ s→−∞<br />
If the interaction V (t) is a bounded operator we may write the unitary<br />
operator in terms of the Dyson series 3<br />
Ũ(t, s) = 1 +<br />
∞<br />
(−i) n<br />
n=1<br />
Claim. t<br />
s<br />
s<br />
s<br />
t<br />
s<br />
dt1<br />
t1<br />
s<br />
tn−1<br />
dt2 · · · dtn<br />
s<br />
˜ V (t1) · · · ˜ V (tn). (3.2)<br />
tn−1<br />
dt1 . . . dtn =<br />
s<br />
1<br />
(t − s)n<br />
n!<br />
Proof. It is obvious for n = 1.<br />
Assume it is right for n = k − 1 then<br />
t tk−1<br />
t<br />
t1 1<br />
dt1 · · · dtk = dt1<br />
(k − 1)!<br />
=<br />
This proves the claim by induction.<br />
s<br />
1<br />
(k − 1)!<br />
t<br />
1<br />
=<br />
(k − 1)!<br />
= 1<br />
(t − s)k<br />
k!<br />
s<br />
t−s<br />
0<br />
s<br />
k−1 dτ<br />
(t1 − s) k−1 dt1<br />
h k−1 dt1<br />
It follows that the sum (3.2) converges in the operator-norm since each<br />
term is bounded by<br />
t<br />
s<br />
dt1 . . .<br />
tn−1<br />
s<br />
dtn ˜ V (t1) . . . ˜ V (tn) = 1<br />
n! (t − s)n ˜ V ) n<br />
If the the interaction decreases for large times such that<br />
then the S-matrix is given by<br />
S =<br />
∞<br />
(−i) n<br />
n=0<br />
+∞<br />
−∞<br />
+∞<br />
−∞<br />
dt1<br />
which is also norm-convergent.<br />
3 See [8] page 326.<br />
t1<br />
dsV (s) < ∞,<br />
dt2 · · ·<br />
−∞<br />
21<br />
tn−1<br />
−∞<br />
dtn ˜ V (t1) · · · ˜ V (tn),
Chapter 4<br />
The Mathematical Setting of<br />
QFT<br />
4.1 The Fock Space<br />
Let H n = H ⊗ · · · ⊗ H be the tensor product space of n single-particle<br />
Hilbert spaces H. Remembering that identical particles must obey either<br />
Bose or Fermi statistics 1 , we symmetrize the states φn ∈ H n<br />
S + n φn = 1<br />
n!<br />
<br />
π<br />
φn(xπ1 , . . . , xπn)<br />
in the case of bosons and antisymmetrize<br />
<br />
(−1) π (xπ1 , . . . , xπn)<br />
S − n φn = 1<br />
n!<br />
π<br />
in the case of fermions. We then form a Hilbert space from the direct sum<br />
of tensor products.<br />
Definition 4.1. The Fock space F is defined as<br />
F ± = ⊕ ∞ n=0S ± n H n<br />
with H0 := {λ|0〉|λ ∈ C} where |0〉 is called the vacuum.<br />
It is easy to prove that F is indeed a Hilbert space. It follows from the<br />
continuity of the inner products. We simply state it here.<br />
Proposition 4.2. The Fock space F is a Hilbert space with the inner<br />
product<br />
∞<br />
〈Φ, Ψ〉 = 〈Φn, Ψn〉n,<br />
n=0<br />
where 〈·, ·〉n is the inner product on Hn.<br />
1 See [8] chapter 6<br />
22
We site Corollary 1.3 of [6] here 2 .<br />
Corollary 4.3. Any (bounded) operator A ∈ F can be expressed as<br />
A = A (−) + A (+)<br />
where A (−) contain annihilation and A (+) creation operators, only.<br />
Definition 4.4. We say that an operator A ∈ F is normally ordered if all<br />
annihilation operators are placed to the right of the creation operators.<br />
The normal ordering of A is expressed by : A :.<br />
In the light of Corollary 4.3 we may define.<br />
Definition 4.5. A contraction between two field operators A and B is<br />
defined by<br />
| AB | := [A (−) , B (+) ]±,<br />
where [·, ·]± is the commutator in the case of Bose fields and<br />
anti-commutator in the case of Fermi fields.<br />
The following theorem tells us how to order field operators normally. We<br />
state it without its proof. 3<br />
Theorem 4.6 (Theorem of Wick). A product of n field operators can be<br />
normally ordered as follows:<br />
A1A2 · · · An =: A1A2 · · · An : + | A1A1 | · · · An + permutations<br />
+ : | A1 | A2 · · · Aj | . . .An : + · · ·<br />
+ : | A1A2 | | A3A4 | · · · : +permutations,<br />
where the sum contains all normal products with all possible pairings of<br />
contractions.<br />
FIXME: think about proving it<br />
We assume the contractions are complex numbers, which therefore can be<br />
taken out of the normal products. Though we have to remember to take<br />
the sign into account in the case of Fermi operators.<br />
In QED there exists only three contractions.<br />
| ψa(x)ψb(y) | := {ψ (−)<br />
a (x), ψ (+)<br />
b (y)} =: 1<br />
i S(+)<br />
ab<br />
| ψa(x)ψb(y) | := {ψ (−)<br />
a (x), ψ (+)<br />
b (y)} =: 1<br />
i S(−)<br />
ba<br />
(x − y) (4.1)<br />
(y − x) (4.2)<br />
| Aµ(x)Aν(y) | := [A (−)<br />
µ (x), A (+)<br />
ν (y)] = gµνiD (+)<br />
0 (x − y), (4.3)<br />
2 I have taken the liberty of reformulating the corollary such that its meaning will be<br />
more apparent for my application. I have not altered the meaning of it, though.<br />
3 See [6]<br />
23
where ψ is the time-dependent Dirac field 4 and ψ = ψ (+) γ 0 is the Dirac<br />
adjoint. The gamma matrices are<br />
γ µ <br />
0 σµ<br />
=<br />
σµ 0<br />
with σµ given by (2.6) and (2.7). We here regard (4.1), (4.2) and (4.3) as<br />
definitions of the operators S = S (−) + S (+) and D (+)<br />
0 , with Aµ given by [6]<br />
page 148.<br />
4 Se [6] page 83<br />
24<br />
<br />
,
4.2 The Wightman Axioms<br />
We will treat a quantum field theory as a tuple<br />
(H, U, φ, D, |0〉),<br />
where H is a separable Hilbert space, U a unitary representation of the<br />
proper Poincaré group P ↑<br />
+ , φ are field operators, D a dense subspace of H<br />
and |0〉 the vacuum state.<br />
We expect the tuple to satisfy the Wightman axioms listed on the next<br />
page. 5<br />
FIXME: make sure the syntax is right<br />
We will not go in to the deeper meaning of the axioms. Instead we will<br />
note the most important properties for our use. We should note that the<br />
fields φ are assumed to be operator valued distributions, which are<br />
well-defined on a dense domain D ⊂ H.<br />
FIXME: note that D0 is contained in axiom 4<br />
In the construction of the scattering matrix, however, it suffices to regard<br />
the dense domain D0 ⊂ D ⊂ H given by,<br />
D0 = span{φ(g1) · · · φ(gn)|0〉|g1, . . . , gn ∈ S, n ∈ N}.<br />
5 We cite them from [9] page 103 − 104.<br />
25
The Wightman Axioms.<br />
Axiom 1 (quantum field) The operators φ1(f), . . . , φn(f) are given<br />
for each C ∞ -function f with compact support on the Minkowski space<br />
R 4 . Each φj(f) and its Hermitian conjugate operator φj(f) ∗ are defined<br />
at least on a common dense linear subset D of the Hilbert space H and<br />
D satisfies<br />
φj(f)D ⊂ D, φj(f) ∗ D ⊂ D,<br />
for any f and j = 1, . . . , n. For any Φ, Ψ ∈ D,<br />
is a complex valued distribution.<br />
f ↦→ (Φ, φj(f)Ψ),<br />
Axiom 2 (relativistic symmetry) On H there exists a unitary repre-<br />
sentation U(a, A) of ˜ P ↑<br />
+ (a ∈ R4 , A ∈ SL(2, C)), satisfying U(a, A)D =<br />
D (invariance of the common domain of the fields).<br />
U(a, A)φj(f)U(a, A) ∗ = S(A −1 )jkφk(f (a,A))<br />
f (a,A)(x) = f(Λ(A) −1 (x − a)),<br />
where the matrix (S(A)j,k) is an n-dimensional representation of A ∈<br />
SL(2C).<br />
Axiom 3 (local commutativity) If the support of f and g is spacelike<br />
separated, then for any vector Φ ∈ D<br />
[φj(f) (∗) , φk(g) (∗) ]±Φ = 0,<br />
where (∗) indicates that the equation holds for any choice.<br />
Axiom 4 (vacuum state) There exists a vector |0〉 in D satisfying the<br />
following conditions:<br />
(i) (Ua, A)|0〉 = |0〉 (invariance).<br />
(ii) The set of all vectors obtained by acting an arbitrary polynomial<br />
P of the fields on |0〉 is dense in H (cyclicity).<br />
(iii) The Spectrum of the translation group U(a, 1) on |0〉 ⊥ is contained<br />
in<br />
V m = {p|(p, p) ≥ m 2 , p 0 > 0} (m > 0)<br />
(Spectrum condition).<br />
26
Chapter 5<br />
The Method of Epstein and<br />
Glaser<br />
5.1 Introduction<br />
We want to express the S-matrix as a formal power series in R[[λ]], where<br />
λ is the coupling constant λ = e, which is the unit of charge and R is the<br />
ring<br />
R = {Υ : D0 → D0|Υ is linear}.<br />
We note that we are not concerned about the convergence of the series and<br />
it might not have any convergence radius at all.<br />
We begin from the expression.<br />
S(g) = 1 +<br />
∞<br />
n=1<br />
=: 1 + T<br />
<br />
n 1<br />
λ<br />
n!<br />
d 4 x1 · · · d 4 xnTn(x1, . . . , xn)g(x1) · · · g(xn)<br />
We shall usually omit the λ in the notation. Further may assume that<br />
Tn(x1, . . . , xn) is symmetric in x1, . . . xn. Otherwise we can always<br />
symmetrize it. This is easy to see in 2 dimensions where we given<br />
T2(x1, x2) can choose the symmetric map defined by<br />
(T2(x1, x2) + T2(x2, x1))/2. The same way we can choose a symmetric map<br />
in n-dimensions as<br />
Tn(x1, x2, . . . , xn) + Tn(x2, x1, x3, . . . , xn) + · · · + Tn(xn, xn−1, . . . , x1)<br />
.<br />
n!<br />
Since Tn(x1, . . . , xn) is symmetric we may occasionally use the short hand<br />
notation Tn(x) where x = {xj ∈ M|j = 1, . . . , n} is disordered.<br />
Along with the Wightman axioms we have to assume that each term of the<br />
S-matrix is well-defined.<br />
27
Axiom 0 (Well-definedness) Let g = g1, . . . , gn ∈ S(R n ), then<br />
〈Tn, (g1 ⊗ · · · ⊗ gn)〉 : D0 → D0,<br />
is a well-defined operator-valued distribution with<br />
<br />
〈Tn, (g1 ⊗ · · · ⊗ gn)〉 = dx1 · · · dxnTn(x1, . . . , xn)g(x1) · · · g(xn).<br />
Note that the S-matrix and Tn are operator-valued distributions, in the<br />
sense that if φ ∈ D0, then<br />
∞<br />
<br />
1<br />
S(g)φ = 1 +<br />
n!<br />
n=1<br />
∞<br />
<br />
1<br />
= φ +<br />
n!<br />
n=1<br />
d 4 x1 . . . d 4 <br />
xnTn(x1, . . . , xn)g(x1) . . . g(xn) φ<br />
d 4 x1 · · · d 4 xnTn(x1, . . . , xn)(φ(g(x1)) · · · φ(g(xn))).<br />
We can express the inverse of S(g) by a similar perturbation series as<br />
follows<br />
S(g) −1 ∞<br />
<br />
1<br />
= 1 + d<br />
n!<br />
n=1<br />
4 x1 . . . d 4 xn ˜ Tn(x1, . . . , xn)g(x1) . . . g(xn)<br />
= (1 + T ) −1 ∞<br />
= 1 + (−T ) r ,<br />
using Theorem 1.7. From<br />
=<br />
∞<br />
n=1<br />
∞<br />
r=1<br />
we see that<br />
n 1<br />
λ<br />
n!<br />
<br />
−<br />
<br />
∞<br />
n=1<br />
r=1<br />
d 4 x1 . . . d 4 xn ˜ Tn(x1, . . . , xn)g(x1) . . . g(xn) =<br />
1<br />
n! λn<br />
<br />
˜Tn(X) =<br />
∞<br />
(−T ) r<br />
r=1<br />
d 4 x1 . . . d 4 r xnTn(x1, . . . , xn)g(x1) · · · g(xn) ,<br />
n<br />
(−1)<br />
r=1<br />
r <br />
where Pr is the set of all partitions of X<br />
Pr<br />
Tn1 (X1) . . . Tnr(Xr), (5.1)<br />
X = X1 ∪ . . . ∪ Xr , |X| = n , Xj = ∅ , |Xj| = nj.<br />
28
Further,<br />
1 = S(g)S(g) −1<br />
=<br />
<br />
1 +<br />
= 1 +<br />
∞<br />
n1=1<br />
<br />
× 1 +<br />
∞<br />
λ n<br />
n=1<br />
<br />
n1 1<br />
λ<br />
n1!<br />
∞<br />
n2=1<br />
n2 1<br />
λ<br />
n1!<br />
<br />
n1+n2=n<br />
d 4 x1 · · · d 4 xn1Tn1 (x1, . . . , xn1 )g(x1) · · · g(xn1 )<br />
<br />
<br />
1<br />
n!<br />
where of course |X| = n1 and T0(∅) = 1 = ˜ T0(∅).<br />
By (5.2)<br />
0 =<br />
∞<br />
λ n<br />
n=1<br />
<br />
n1+n2=n<br />
<br />
<br />
1<br />
n!<br />
d 4 ˜x1 · · · d 4 ˜xn2 ˜ Tn2 (˜x1, . . . , ˜xn2 )g(˜x1) · · · g(˜xn2 )<br />
<br />
d 4 x1 · · · d 4 xnTn1 (X) ˜ Tn2 ( ˜ X)g(x1) · · · g(xn),<br />
d 4 x1 · · · d 4 xnTn1 (X) ˜ Tn2 ( ˜ X)g(x1) · · · g(xn)<br />
Then by the definition of the metric on formal power series each power<br />
must vanish. That is,<br />
<br />
P 0 2<br />
(5.2)<br />
Tn1 (X) ˜ Tn−n1 ( ˜ X)g(x1) · · · g(xn) = 0, (5.3)<br />
where P 0 2 is the set of all partitions {x1, . . . , xn} = X ∪ ˜ X.<br />
In constructing the scattering matrix, we should be aware of the properties<br />
we think it should have. We list such expected properties on the next page.<br />
29
Expected properties of the S-matrix.<br />
1 Unitarity i.e. S(g) −1 = S(g) ∗ . That is<br />
˜Tn(x1, . . . , xn) = Tn(x1, . . . , xn) ∗ .<br />
Fixme: this is not enough according to page 162 probably not important for<br />
this project.<br />
2 Translational invariance. Let U(a, 1) be the unitary translation operator<br />
in the Fock space F, that is<br />
Then we require<br />
(U(a, 1)Φ)j(x) = Φj(x1 + a, . . . , xj + a).<br />
U(a, 1)S(g)U(a, 1) −1 = S(ga), where ga(x) = g(x − a).<br />
Hence for the Tn’s (and of course the ˜ Tn’s as well) we get:<br />
U(a, 1)Tn(x1, . . . , xn)U(a, 1) −1 = Tn(x1 + a, . . . , xn + a).<br />
3 Lorentz covariance Letting U(0, Λ) be the representation of L ↑<br />
+ .<br />
U(0, Λ)S(g)U(0, Λ) −1 = S(gΛ), where gΛ = g(Λ −1 x).<br />
Note that 2. and 3. together form a condition of Poincaré invariance.<br />
4 Causality. Suppose there exists a reference frame in which the<br />
test-functions g1 and g2 have disjoint supports in time. Assuming<br />
supp g1 < supp g2. That is, for some r ∈ R,<br />
supp g1 ⊂ {x ∈ M|x 0 ∈ (−∞, r)} and supp g2 ⊂ {x ∈ M|x 0 ∈ (r, ∞)}.<br />
We require that<br />
S(g1 + g2) = S(g2)S(g1). (5.4)<br />
This is a statement about the fact that what happens at time t < s is not<br />
influenced by what happens at some later time t > s.<br />
30
Of these properties, causality plays a key role in the method of Epstein<br />
and Glaser. Equation (5.4) leads to the condition stated in the following<br />
theorem.<br />
Theorem 5.1. If {x1, . . . , xm} > {xm+1, . . . , xn}, then<br />
Tn(x1, . . . , xn) = Tm(x1, . . . , xm)Tn−m(xm+1, . . . , xn),<br />
Before proving this result we first note that if (5.4) is satisfied,<br />
S(g1 + g2) =<br />
=<br />
∞<br />
n=0<br />
<br />
n 1<br />
λ<br />
n!<br />
d 4 x1 · · · d 4 xnTn(x1, . . . , xn)<br />
× (g1(x1) + g2(x1)) · · · (g1(xn) + g2(xn))<br />
∞ n<br />
λ n<br />
<br />
1<br />
m!(n − m)!<br />
n=0 m=0<br />
d 4 x1 · · · d 4 xnTn(x1, . . . , xn)<br />
× (g2(x1) · · · g2(xm)) · · · (g1(xm+1) · · · g1(xn)),<br />
as there are 2n terms from the product of the test-functions and<br />
ways to pick g2 m times. On the other hand<br />
S(g2)S(g1) =<br />
=<br />
=<br />
∞<br />
m=0<br />
<br />
m 1<br />
λ<br />
m!<br />
d 4 x1 · · · d 4 xmTm(x1, . . . , xm)<br />
× g2(x1) · · · g2(xm)<br />
∞<br />
<br />
k 1<br />
× λ<br />
k!<br />
k=0<br />
∞<br />
m=0 k=0<br />
∞<br />
n=0 m=0<br />
d 4 ˜x1 · · · d 4 ˜xkTk(˜x1, . . . , ˜xk)<br />
× g1(˜x1) · · · g1(˜xk)<br />
∞<br />
<br />
m+k 1<br />
λ<br />
m!k!<br />
d 4 x1 · · · d 4 xmd 4 ˜x1 · · · d 4 ˜xk<br />
× Tm(x1, . . . , xm)Tk(˜x1, . . . , ˜xk)<br />
n!<br />
m!(n−m)!<br />
× g2(x1) · · · g2(x1) · · · g2(xm)g1(˜x1) · · · g1(˜xk)<br />
n<br />
λ n<br />
<br />
1<br />
dx1 · · · dxnTm(x1, . . . , xm)<br />
m!(n − m)!<br />
× Tn−m(xm+1, . . . , xn)g2(x1) · · · g2(xm)<br />
× g1(xm+1) · · · g1(xn).<br />
From the definition of the metric on formal power series we know that<br />
31
equal powers must have equal coefficients. We conclude that for every m, n<br />
<br />
<br />
=<br />
d 4 x1 · · · d 4 xnTn(x1, . . . , xn)g2(x1) · · · g2(xm)g1(xm+1) · · · g1(xn)<br />
d 4 x1 · · · d 4 xnTm(x1, . . . , xm)Tn−m(xm+1, . . . , xn)<br />
× g2(x1) · · · g2(xm)g1(xm+1) · · · g1(xn) (5.5)<br />
Further note that from the calculus of the tensor product<br />
〈T2, (f1 + f2) ⊗ (f1 + f2)〉 = 〈T2, f1 ⊗ f1〉 + 〈T2, f1 ⊗ f2〉<br />
+ 〈T2, f2 ⊗ f1〉 + 〈T2, f2 ⊗ f2〉,<br />
FIXME: maybe better to write it out in n dimensions<br />
and since Tm is symmetric<br />
〈Tm, f1 ⊗ f2〉 = 1<br />
2 (〈Tm, (f1 + f2) ⊗ (f1 + f2)〉 − 〈Tm, f1 ⊗ f1〉<br />
− 〈Tm, f2 ⊗ f2〉). (5.6)<br />
Now for transparency and to get a good idea of how to prove Theorem 5.1<br />
let us prove it for n = 3 and m = 2. That is,<br />
Claim. 〈T3, f1 ⊗ f2 ⊗ f3〉 = 〈T2, f1 ⊗ f2〉〈T1, f3〉.<br />
Proof. By 5.6<br />
Hence by 5.5<br />
〈T2, f1 ⊗ f2〉 = 1<br />
〈T2, (f1 + f2) ⊗ (f1 + f2)〉 − 〈T2, f1 ⊗ f1〉<br />
2<br />
− 〈T2, f2 ⊗ f2〉 <br />
〈T2, f1 ⊗ f2〉〈T1, f3〉<br />
= 1<br />
〈T3, (f1 + f2) ⊗ (f1 + f2) ⊗ f3〉 − 〈T3, f1 ⊗ f1 ⊗ f3〉<br />
2<br />
− 〈T3, f2 ⊗ f2 ⊗ f3〉 <br />
= 1<br />
〈T3, f1 ⊗ f1 ⊗ f3〉 + 〈T3, f2 ⊗ f2 ⊗ f3〉 + 2〈T3, f1 ⊗ f2 ⊗ f3〉<br />
2<br />
− 〈T3, f1 ⊗ f1 ⊗ f3〉 − 〈T3, f2 ⊗ f2 ⊗ f3〉 <br />
= 〈T3, f1 ⊗ f2 ⊗ f3〉.<br />
Which proves the claim.<br />
Now before proving this for arbitrary n and m we need the following<br />
lemma.<br />
32
Lemma 5.2. The product (x1 + x2 + . . . + xn) m has<br />
n−1 <br />
m m − k1 m − i<br />
ki<br />
· · · ·<br />
,<br />
k1<br />
terms on the form x k1<br />
1<br />
· · · xkn<br />
n .<br />
k2<br />
Proof. We note that the number of ways to pick x1, k1 times from m<br />
possibilities<br />
<br />
is given by the binomial coefficient. That is, there must be<br />
m<br />
k1<br />
terms on the form x k1<br />
1 (x2 + · · · + xn) m−k1 . Now there are m − k1<br />
multiplications left and thus m−k1 ways to pick x2, k2 times. Hence there<br />
k2<br />
are m m−k1<br />
k1<br />
terms on the form x k1 k2<br />
1 xk2 2 (x3 + · · · xn) m−k1−k2 . This way we<br />
end up with<br />
n−1 <br />
m m − k1 m − i<br />
ki<br />
· · · ·<br />
,<br />
k1<br />
terms on the form x k1<br />
1<br />
· · · xkn<br />
n .<br />
Proof of Theorem 5.1. First note that by the lemma<br />
k2<br />
〈Tm, (f1 + f2 + · · · + fm) ⊗ · · · ⊗ (f1 + f2 + · · · + fm)〉<br />
=<br />
<br />
n−1 <br />
m m − k1 m − i<br />
ki<br />
· · · ·<br />
k1+···+km=m<br />
k1<br />
since Tm is symmetric. Hence<br />
〈Tm, f1 + f2 + · · · + fm〉<br />
k2<br />
kn<br />
kn<br />
× 〈Tm, f k1<br />
1<br />
kn<br />
⊗ · · · ⊗ f km<br />
m 〉,<br />
= 〈Tm, (f1 + f2 + · · · + fm) ⊗ . . . ⊗ (f1 + f2 + · · · + fm)〉<br />
−<br />
<br />
n−1 <br />
m m − k1 m − i<br />
ki<br />
· · · ·<br />
〈Tm, f k1<br />
1<br />
k1+···+km=m<br />
(k1,...,km)=(1,...,1)<br />
Hence by 5.5<br />
k1<br />
〈Tm, f1 ⊗ · · · ⊗ fm〉〈Tn−m, fm+1 ⊗ · · · ⊗ fn〉<br />
k2<br />
kn<br />
⊗ · · · ⊗ f km<br />
m 〉.<br />
= 〈Tn, (f1 + f2 + · · · + fm) ⊗ . . . ⊗ (f1 + f2 + · · · + fm) ⊗ (fm+1 ⊗ · · · ⊗ fn)〉<br />
−<br />
<br />
n−1 <br />
m m − k1 m − i<br />
ki<br />
· · · ·<br />
k1+···+km=m<br />
(k1,...,km)=(1,...,1)<br />
× 〈Tn, f k1<br />
1<br />
k1<br />
k2<br />
⊗ f km<br />
m ⊗ (fm+1 ⊗ · · · ⊗ fn)〉<br />
= 〈Tn, f1 ⊗ f2 ⊗ · · · ⊗ fm ⊗ fm+1 ⊗ · · · ⊗ fn〉<br />
as wanted.<br />
33<br />
kn
We see that the Tn’s are time-ordered products (therefore the T ).<br />
Similarly the causality condition for S −1 (g) is that if supp g1 < supp g2,<br />
then<br />
S(g1 + g2) −1 = S(g1) −1 S(g2) −1 .<br />
This again implies that<br />
˜Tn(x1, . . . , xn) = ˜ Tm(x1, . . . , xm) ˜ Tn−m(xm+1, . . . , xn),<br />
if {x1, . . . , xm} < {xm+1, . . . , xn}.<br />
We are now ready to sketch the game plan for the inductive construction<br />
of the time ordered products Tn.<br />
34
The Inductive Construction of Tn(x1, . . . , xn).<br />
1. Assume Tm(x1, . . . , xm) for 1 ≤ m ≤ n − 1 are known.<br />
2. Construct advanced and retarded distributions as follows<br />
A ′ (x1, . . . , xn) = <br />
˜Tn1 (X)Tn−n1 (Y, xn) (5.7)<br />
P2<br />
R ′ (x1, . . . , xn) = <br />
P2<br />
Tn−n1 (Y, xn) ˜ Tn1 (X) (5.8)<br />
where P2 is the set of all partitions {x1, . . . , xn−1} = X ∪ Y , X = ∅<br />
and n1 = |X| ≥ 1.<br />
3. Include the empty set ∅ as follows<br />
A(x1, . . . , xn) = <br />
˜Tn1 (X)Tn−n1 (Y, xn)<br />
P 0 2<br />
= A ′ n(x1, . . . , xn) + Tn(x1, . . . , xn) (5.9)<br />
R(x1, . . . , xn) = <br />
Tn−n1 (Y, xn) ˜ Tn1 (X)<br />
4. The difference is<br />
5. Now<br />
P 0 2<br />
= R ′ n(x1, . . . , xn) + Tn(x1, . . . , xn). (5.10)<br />
Dn = R ′ n − A ′ n = Rn − An. (5.11)<br />
Tn = Rn − R ′ n = An − A ′ n. (5.12)<br />
We thus need to determine either Rn or An. This will be done by<br />
investigating the support properties of the distributions.<br />
35
5.2 Example - In the Hilbert Space Setting of<br />
Quantum Mechanics<br />
Before going on we illustrate the significance of the table by applying it to<br />
the more transparent theories of Hilbert spaces and quantum mechanics.<br />
In quantum mechanics we substitute distributions by functions and replace<br />
x by t.<br />
Now, assuming T1(t) is given we want to construct T2(t). We follow the<br />
rules given from the box on the previous page and construct the advanced<br />
and retarded functions.<br />
and<br />
A ′ 2(t1, t2) = ˜ T1(t1)T1(t2) = −T1(t1)T1(t2) (5.13)<br />
R ′ 2(t1, t2) = T1(t2) ˜ T1(t1) = −T1(t2)T1(t1), (5.14)<br />
where the last inequalities follow from (5.1). Then<br />
and<br />
A2(t1, t2) = A ′ 2(t1, t2) + T2(t1, t2) (5.15)<br />
R2(t1, t2) = R ′ 2(t1, t2) + T2(t2, t1)<br />
From Theorem 5.1 we know that if ti > tj, then<br />
T2(ti, tj) = T1(t1)T1(tj).<br />
Hence A2 vanishes for t1 > t2 and R2 vanishes for t1 < t2.<br />
Now<br />
D2 = R2 − A2 = R ′ 2 − A ′ 2<br />
is known from equations (5.13) and (5.14). The T2 cancel out since they<br />
are symmetric.<br />
Finally, R2 and A2 can be determined by their support properties. Let<br />
Θ(t) be the Heaviside function<br />
Then<br />
Θ(x) =<br />
<br />
1, t ≥ 0,<br />
0, t < 0.<br />
A2(t1, t2) = Θ(t2 − t1)D2(t1, t2)<br />
= Θ(t2 − t1)(T1(t1)T1(t2) − T1(t2)T1(t1)).<br />
Hence A2 is uniquely determined up to its value at t1 = t2.<br />
36
Now from equation (5.15)<br />
T2(t1, t2) = A2(t1, t2) − A ′ 2(t1, t2)<br />
= Θ(t2 − t1)(T1(t1)T1(t2) − T1(t2)T1(t1)) + T1(t1)T1(t2)<br />
= Θ(t1 − t2)T1(t1)T1(t2) − Θ(t2 − t1)T1(t2)T1(t1)<br />
=: T {T1(t1)T1(t2)}.<br />
FIXME: draw parallel to S-matrix with usage of Schrödinger equation.<br />
37
5.3 Splitting of Distributions<br />
In our previous example we saw that the splitting of Dn into advanced and<br />
retarded parts could be done by use of the Heaviside function. This was<br />
because we were dealing with operators in Hilbert space. But in our<br />
setting we are considering distributions which cannot simply be multiplied<br />
with discontinuous step-functions. In this section we consider the<br />
properties of the advanced and retarded distributions.<br />
First a technicality.<br />
Theorem 5.3. Let Y = P ∪ Q, P = ∅, P ∩ Q = ∅, |Y | = n1 ≤ n − 1 and<br />
x /∈ Y .<br />
If {Q, x} > P ,|Q| = n2, then we have<br />
If {Q, x} < P , then we have<br />
R ′ n1+1(Y, x) = −Tn2+1(Q, x)Tn1−n2 (P ). (5.16)<br />
A ′ n1+1(Y, x) = −Tn1−n2 (P )Tn2+1(Q, x). (5.17)<br />
Proof. Equation (5.16) is proved in [6]. Therefore we only prove equation<br />
(5.17).<br />
A ′ n1+1(Y, x) = <br />
˜Tn3 (X)Tn−n3 (Y ′ , x),<br />
where P2 is the set of all partitions Y = X ∪ Y ′ such that X = ∅.<br />
Now let<br />
and<br />
Now since<br />
causality implies<br />
P2<br />
Y ′ = Y1 ∪ Y2, where Y1 = Y ′ ∩ P and Y2 = Y ′ ∩ Q<br />
X = X1 ∪ X2, where X1 = X ′ ∩ P and X2 = X ∩ Q<br />
Q ∪ {x} < P, Y2 < Y1, X2 < X1 and Y1 > Y2 ∪ {x},<br />
A ′ n1+1(Y, x) = <br />
˜Tn3 (X)Tn−n3 (Y ′ , x)<br />
P2<br />
= <br />
˜T (X2) ˜ T (X1)T (Y1)T (Y2, x), (5.18)<br />
P 0 4<br />
subscripts are omitted for simplicity. P 0 4<br />
form<br />
is the set of all partitions of the<br />
P 0 4 : P = X1 ∪ Y1, Q = X2 ∪ Y2 and X1 ∪ X2 = ∅.<br />
38
However, for x2 = ∅ we can have X1 = ∅. From Equation (5.3) it follows<br />
that <br />
˜T (X1)T (Y1) = 0. (5.19)<br />
P 0 2<br />
As a consequence, in (5.18) only terms with X2 = ∅ and hence Y2 = Q<br />
remain. That is, remembering that T0(∅) = 1 = ˜ T0(∅),<br />
A ′ <br />
<br />
n1+1(Y, x) = ˜T (X1)T (Y1) T (Q, x),<br />
P2<br />
where P2 is set set of all partitions such that P = X1 ∪ Y1 and X1 = ∅.<br />
Including ∅ would give 0 according to (5.19), hence<br />
A ′ <br />
n1+1(Y, x) = ˜T (X1)T (Y1) − ˜ <br />
T (∅)T (P ) T (Q, x)<br />
as wanted.<br />
P 0 2<br />
Fixme: investigate the meaning of this.<br />
= (0 − T (P ))T (Q, x) = −T (P )T (Q, x),<br />
Corollary 5.4 (Support property 1). The supports of A and R are<br />
respectively advanced and retarded, that is<br />
and<br />
Proof. Equation (5.9) gives us<br />
An1+1(Y, x) = A ′ n1+1 + Tn(Y, x)<br />
supp An1+1(Y, x) ⊂ Γ − n1+1 (x)<br />
supp Rn1+1(Y, x) ⊂ Γ + n1+1 (x)<br />
= −Tn1−n2 (P )Tn2+1(Q, x) + Tn1+1(P ∪ Q, x)<br />
= −Tn1−n2 (P )Tn2+1(Q, x) + Tn1−n2 (P )Tn2+1(Q, x) = 0.<br />
2. follows the same way from (5.10).<br />
The corollary explains why A and R are called advanced and retarded<br />
distributions. An+1(X, x) vanishes if there exists a point x ′ ∈ X with<br />
x ′ > x, i.e if P exists. On the other hand Rn+1(X, x) vanishes if there<br />
exists a point x ′ ∈ X with x ′ < x.<br />
Theorem 5.5 (Support property 2). If n ≥ 3 then<br />
supp Dn(x1, . . . , xn−1, xn) ⊆ Γ − n (xn) ∪ Γ + n (xn). (5.20)<br />
39
A proof of this is found in [6] page 167. For n ≤ 2 this property has to be<br />
verified explicitly.<br />
It is here the difficult part of the method of Epstein and Glaser described<br />
in [6] lies, that is in decomposing<br />
such that<br />
Dn(x1, . . . , xn) = Rn(x1, . . . , xn) − An(x1, . . . , xn)<br />
supp Rn ⊆ Γ + n−1 (xn) and supp An ⊆ Γ − n−1 (xn).<br />
40
Chapter 6<br />
Regularly Varying Functions<br />
Before going on in the splitting process we need to consider some<br />
properties of a class of functions which will become important.<br />
Definition 6.1. A positive function ρ(t) is called regularly varying at<br />
t = 0, if it is measurable in (0, t0] for some t0 > 0, and there exists ω ∈ R<br />
such that the limit<br />
ρ(at)<br />
lim = aω<br />
t→0 ρ(t)<br />
exists for all a > 0. In this case ω is called the order of ρ. If ω = 0 the<br />
function is said to be slowly varying.<br />
(6.1)<br />
Note that a regularly varying function always can be reduced to be slowly<br />
varying. That is, let<br />
ρ(t) = t ω r(t). (6.2)<br />
Then (6.1) gives<br />
r(at)<br />
lim = 1, (6.3)<br />
t→0 r(t)<br />
for all a > 0.<br />
The following representation theorem will be important for our use.<br />
Theorem 6.2. Let r : (0, t0] → R+ be a slowly varying function. Then<br />
there exists some t1 ∈ (0, t0) such that<br />
t1 h(s)<br />
r(t) = exp η(t) +<br />
t s ds<br />
<br />
holds for all 0 < t < t1. Here η is bounded and measurable on (0, t1], and<br />
η → c for t → 0 (|c| < ∞). Further h(t) is continuous on (0, t1] and<br />
h(t) → 0 for t → 0.<br />
Before proving the theorem we will need some lemmas.<br />
41
Lemma 6.3. Let f : [x0, ∞) → R be a measurable function such that<br />
f(x + k) − f(x) → 0 for x → ∞ (6.4)<br />
for all k ∈ R. 1 Let I = [a, b] ⊂ R be a closed interval then<br />
Proof. Let I = [0, 1]. Assume<br />
sup |f(x + k) − f(x)| → 0 for x → ∞.<br />
k∈I<br />
sup |f(x + k) − f(x)| 0<br />
k∈I<br />
for x → ∞. (6.5)<br />
Then there exists an ɛ > 0 and a sequence (xn, kn) where xn → ∞ and<br />
kn ∈ I such that |f(xn + kn) − f(xn)| ≥ ɛ for all n ∈ N.<br />
Now consider the measurable sets<br />
Fixme: are they?<br />
and<br />
Un = {s ∈ [0, 2]||f(xm + s) − f(xm)| < ɛ/2 for all m ≥ n}<br />
Vn = {t ∈ [0, 2]||f(xm + km + t) − f(xm + km)| < ɛ/2 for all m ≥ n}.<br />
By (6.4) these sets are clearly monotonically increasing towards [0, 2]<br />
hence, letting m denote the Lebesgue measure, there exists an N ∈ N such<br />
that m(Un) > 3/2 and m(Vn) > 3/2 for n, m ≥ N.<br />
Consider the translation V ′ N = VN + kN, then since m(V ′ N ) > 3/2 and<br />
UN, V ′ N ⊂ [0, 3] the two sets must intersect in some point k. Now since<br />
k ∈ UN<br />
|f(xN + k) − f(xN)| < ɛ/2<br />
and since k − kN ∈ VN<br />
|f(xN + k) − f(xN + kN)| = |f(xN + kN + k − kN) − f(xN + kN)| < ɛ/2.<br />
Hence<br />
|f(xN + kn) − f(xN)| = |f(xN + k) − f(xN) − (f(xN + k) − f(xN + kN))|<br />
≤ |f(xN + k) − f(xN)| + |f(xN + k) − f(xN + kN)|<br />
< ɛ/2 + ɛ/2 = ɛ,<br />
contradicting (6.5). This proves the theorem for I = [a, b] where a = 0 and<br />
b = 1.<br />
1 For k < 0 we may only consider x large enough that f be defined. This is no problem,<br />
though, since we are only concerned about the limit x → ∞.<br />
42
Now we want to use this to show it for arbitrary a, b ∈ R+. Let<br />
Then for a fixed,<br />
g(x) := f((b − a)x).<br />
f(x + k) − f(x) = f(x + a + k − a) − f(x + a) + f(x + a) − f(x)<br />
<br />
x + a k − a<br />
<br />
x + a<br />
<br />
= g + − g + f(x + a) − f(x)<br />
b − a b − a b − a<br />
= g(y + h) − g(y) + f(x + a) − f(x), (6.6)<br />
where y and h are defined in the obvious way. Now since k ∈ [a, b] we see<br />
that h ∈ [0, 1]. Further y → ∞ is equivalent to x → ∞. Now by (6.4),<br />
f(x − a) − f(x) → 0 for x → ∞. And by the proof above for I = [0, 1]<br />
Hence by (6.6)<br />
as claimed.<br />
sup |g(y + h) − g(y)| → 0 for y → ∞.<br />
h∈I<br />
sup<br />
k∈[a,b]<br />
|f(x + k) − f(x)| → 0 for x → ∞.<br />
Lemma 6.4. Let f(x) = ln r(e −x ), where r is the slowly varying function<br />
defined by (6.2). Then f satisfies (6.4) and there exists a constant x0 such<br />
that f is bounded on every interval [a, b] where x0 ≤ a.<br />
FIXME: x0 comes from lemma 6.3 and tells us what the domain of f is.<br />
Proof.<br />
lim (f(x + k) − f(x)) = lim<br />
x→∞ x→∞ (ln r(e−(x+k) ) − ln r(e x ))<br />
r(e<br />
= ln lim<br />
x→∞<br />
−(x+k) )<br />
= ln lim<br />
t→0<br />
= ln 1 = 0,<br />
r(e −x )<br />
since (6.3) holds for r(t). Therefore f satisfies (6.4).<br />
By Lemma 6.3 there exists a such that<br />
for all x ≥ a and for all k ∈ [0, 1].<br />
r(at)<br />
r(t) , substituting t = e−x and a = e −k<br />
|f(x + k) − f(x)| < 1 (6.7)<br />
43
We want to show by induction that for x ∈ [a + m − 1, a + m],<br />
|f(x)| ≤ |f(a)| + m. (6.8)<br />
For m = 1, let x ∈ [a, a + 1] then we may write x = a + k for some<br />
k ∈ [0, 1]. Then<br />
|f(a + k)| = |f(a + k) − f(a) + f(a)| ≤ |f(a + k) − f(a)| + |f(a)|<br />
≤ |f(a)| + 1.<br />
In the (n + 1)th step, let x ∈ [a + n, a + n + 1]. Then there exists a<br />
kn ∈ [0, 1] such that x = a + n + kn and<br />
|f(x)| = |f(a + n + kn)| = |f(a + n + kn) − f(a + n) + f(a + n)|<br />
≤ |f(a + n + kn) − f(a + n)| + |f(a + n)|<br />
< |f(a)| + n + 1<br />
by the induction hypothesis. This concludes the proof by induction.<br />
Clearly if x satisfies (6.8), then it is also satisfied by all y ∈ [a, a + m] and<br />
we have proven that f is bounded on the interval [a, b].<br />
Corollary 6.5. Let f(x) = ln r(e −x ). Then there exists a constant x0 such<br />
that f is integrable on every interval [a, b] where x0 ≤ a.<br />
Proof. This follows directly from the definition of integrability since f is<br />
measurable and bounded on [a, b] by Lemma 6.4.<br />
Lemma 6.6. Let a ≥ x0 from Lemma 6.4. Then f(x) = ln r(e−x ) can be<br />
represented for all x ≥ a as<br />
f(x) = g(x) +<br />
x<br />
a<br />
h(s)ds, (6.9)<br />
where g and h are measurable and bounded on every interval [a, b] and<br />
g(x) → g0 (|g0| ≤ ∞) and h(x) → 0 for x → ∞.<br />
Proof. Let x ≥ a we may write f(x) on the form (6.9) by letting<br />
g(x) :=<br />
x+1<br />
FIXME: maybe write out the above<br />
By (6.4) h(x) → 0 for x → ∞.<br />
x<br />
f(x) − f(s)ds +<br />
a+1<br />
h(x) := f(x + 1) − f(x)<br />
44<br />
a<br />
f(s)ds
By Lemma 6.4 the second integral of g(x) is bounded. Consider the first<br />
integral<br />
x+1<br />
x<br />
f(x) − f(s)ds =<br />
≤<br />
1<br />
0<br />
1<br />
by (6.7). Hence g is bounded. Further<br />
lim k(x) = lim<br />
x→∞ x→∞<br />
=<br />
0<br />
1<br />
0<br />
f(x) − f(x + k)dk, substituting s = x + k<br />
|f(x) − f(x + k)|dk < 1,<br />
1<br />
0<br />
f(x) − f(x + k)dk<br />
lim (f(x) − f(x + k))dk = 0,<br />
x→∞<br />
by (6.4) using the Dominated Convergence Theorem with (6.7). We<br />
conclude that<br />
as claimed.<br />
g(x) →<br />
a+1<br />
a<br />
f(s)ds =: g0 for x → ∞, (6.10)<br />
Lemma 6.7. There exists a∗ such that for all x ≥ a∗ x<br />
f(x) = g ∗ (x) +<br />
a ∗<br />
h ∗ (s)ds<br />
where g ∗ and h ∗ are measurable and bounded on every interval [a ∗ , b] and<br />
g ∗ (x) → g ∗ 0 (|g∗ 0 | ≤ ∞) and h∗ (x) → 0 for x → ∞. In fact, h ∗ is continuous<br />
on every interval [a ∗ , b].<br />
FIXME: reformulate.<br />
Proof. Let f ∗ (x) = x<br />
a h(s)ds. Then by (6.9) and (6.10)<br />
For all k > 0,<br />
f(x) − f ∗ (x) = g(x) → g0 for x → ∞.<br />
f ∗ (x + k) + f ∗ x+k<br />
(x) =<br />
a<br />
x+k<br />
=<br />
=<br />
=<br />
x<br />
x+k<br />
x<br />
k<br />
0<br />
h(s)ds +<br />
h(s)ds<br />
a<br />
x<br />
h(s)ds<br />
f(s + 1) − f(s)ds<br />
f(t + x + 1) − f(t + x)dt,<br />
45
substituting t := s − x. We may write<br />
f(t + x + 1) − f(t + x) = [f(x + t + 1) − f(x)] − [f(x + t) − f(x)].<br />
Both terms on the right hand side converge uniformly in t towards 0 on<br />
[0, k] as x → ∞ by Lemma 6.3. Hence<br />
f ∗ (x + k) − f ∗ k<br />
(x) =<br />
0<br />
f(x + t + 1) − f(x + t)<br />
≤ k sup |f(x + t + 1) − f(x + t)| → 0<br />
t∈[0,k]<br />
for x → ∞. The same way it can be shown to hold for k < 0 as well.<br />
Hence f ∗ satisfies (6.4).<br />
Note that f ∗ is continuous in any interval [a, b], because given ɛ > 0,<br />
choose δ = ɛ/ sup [a,b] |f(s + 1) − f(s)|. Then<br />
x0<br />
|f(x0) − f(x)| = |<br />
≤<br />
x x0<br />
x<br />
f(s + 1) − f(s)ds|<br />
|f(s + 1) − f(s)|ds<br />
≤ |x0 − x| sup<br />
[a,b]<br />
|f(s + 1) − f(s)|<br />
< δ sup |f(s + 1) − f(s)| < ɛ.<br />
[a,b]<br />
Now we want to construct a representation of f ∗ .<br />
First note that since f ∗ is continuous it is measurable.<br />
Thus, by Lemma 6.3 we may conclude that<br />
sup |f<br />
[c,d]<br />
∗ (x + k) − f ∗ (x)| → 0 for x → ∞.<br />
Further, since by Lemma 6.4 f is bounded on every closed interval so is f ∗ ,<br />
because<br />
sup<br />
[c,d]<br />
|f ∗ (x)| = sup |<br />
x<br />
[c,d] a<br />
x<br />
≤ sup<br />
[c,d]<br />
a<br />
f(s + 1) − f(s)ds|<br />
|f(s + 1) − f(s)|ds<br />
≤ (d − a) sup |f(t + 1) − f(t)| < ∞.<br />
[a,d]<br />
Since f ∗ is measurable and bounded on closed intervals, it is integrable.<br />
By Lemma 6.6 f ∗ has the following representation for x ≥ a ∗ ,<br />
FIXME: what is a ∗ ?.<br />
46
where<br />
and<br />
k ∗ (x) =<br />
f ∗ (x) = k ∗ x<br />
(x) +<br />
a∗ h ∗ (s)ds,<br />
x+1<br />
f<br />
x<br />
∗ (x) − f ∗ a∗ +1<br />
(s)ds +<br />
a∗ h ∗ (x) = f ∗ (x + 1) − f ∗ (x).<br />
Note that h∗ is continuous, since f ∗ is continuous.<br />
Now (6) implies<br />
f(x) = g(x) + f ∗ (x) = g(x) + g ∗ x<br />
(x) +<br />
a∗ Finally, let g ∗ (x) := g(x) + k ∗ (x).<br />
f ∗ (s)ds<br />
h ∗ (s)ds<br />
Now we are finally ready to prove the representation theorem.<br />
Proof of Theorem 6.2. Let f(x) := ln r(e−x ). By Lemma 6.7 f has the<br />
following representation<br />
That is<br />
ln(r(e −x )) = g ∗ x<br />
(x) +<br />
a∗ ln r(t) = g ∗ (− ln t) +<br />
− ln t<br />
a ∗<br />
h ∗ (s)ds.<br />
h ∗ (s)ds.<br />
Letting ν(t) := g ∗ (− ln(t)), substituting s := − ln t ′ and h(t ′ ) := h ∗ (− ln t ′ ),<br />
Thus,<br />
where t1 = e −a∗<br />
t<br />
ln r(t) = η(t) −<br />
e−a∗ h(t ′ )<br />
dt′<br />
t ′<br />
t1 h(t<br />
r(t) = exp(η(t) +<br />
t<br />
′ )<br />
t ′ dt′ ),<br />
as claimed.<br />
Note that the representation of 6.2 only holds for t < t1 but since we are<br />
only concerned with the limit t → 0 we will usually not have any problem<br />
using the representation theorem. The following theorem will be important<br />
later when we shall split distributions.<br />
Corollary 6.8. Let ρ be a regularly varying function and ɛ > 0, then<br />
there exist some S0(ɛ) and non-negative constants C and C ′ such that<br />
C ′ t ω+ɛ ≤ ρ(t) ≤ Ct ω−ɛ , for 0 ≤ t ≤ s0(ɛ). (6.11)<br />
47
Proof. Reduce ρ to be slowly varying<br />
ρ(t) = t ω r(t) = t ω t1<br />
exp η(t) +<br />
t<br />
h(s)<br />
s ds<br />
<br />
, for 0 < t < t1,<br />
where we substitute r by the representation given by Theorem 6.2. Since<br />
h(s) → 0 for s → 0 there exists s0(ɛ) ≤ t1 such that |h(s)| ≤ ɛ for s ≤ s0(ɛ).<br />
So if 0 < t ≤ s0(ɛ), then<br />
<br />
<br />
<br />
<br />
s0(ɛ)<br />
t<br />
h(s)<br />
s ds<br />
<br />
<br />
<br />
≤<br />
s0(ɛ)<br />
t<br />
<br />
<br />
<br />
h(s) <br />
<br />
s ds ≤<br />
s0(ɛ)<br />
t<br />
ɛ<br />
s ds = ɛ ln s0(ɛ) − ɛ ln t<br />
<br />
s0(ɛ)<br />
= ɛ ln .<br />
t<br />
Hence<br />
ρ(t) = t ω s0(ɛ) h(s)<br />
exp η(t) +<br />
t s ds<br />
<br />
≤ t ω <br />
s0(ɛ)<br />
exp η(t) + ɛ ln<br />
t<br />
= t ω−ɛ e η(t) s ɛ 0(ɛ)<br />
≤ Ct ω−ɛ ,<br />
where C = s ɛ 0 (ɛ) sup 0
Chapter 7<br />
Splitting of Numerical<br />
Distributions<br />
7.1 The Singular Order of a Distribution<br />
We expect the operator-valued distributions to be expanded in terms of<br />
free fields as follows1 Dn(x1, . . . , xn) = <br />
: <br />
ψ(xj)d k n(x1, . . . , xn) <br />
ψ(xl) :: <br />
A(xm) :,<br />
k<br />
j<br />
where the numerical distributions d k n ∈ S ′ (R 4n ) have the causal support<br />
property (5.20). We have to assume them to be tempered in order to use<br />
the Fourier transformation. Our aim is to split these distributions<br />
where<br />
d k n(x) = rn(x) − an(x),<br />
supp rn ⊂ Γ + n−1 (xn) and supp an ⊂ Γ − n−1 (xn).<br />
We have assumed translational invariance thus we may assume xn = 0 and<br />
thus only consider<br />
d(x) := d(x1, . . . xn−1, 0) ∈ S ′ (R m ), m = 4(n − 1).<br />
In the Hilbert space setting of our earlier example we would do the<br />
splitting as<br />
rn(x) = χn(x)d k n(x), where χn(x) =<br />
n−1 <br />
j=1<br />
l<br />
Θ(x 0 j − x 0 n) =<br />
m<br />
n−1 <br />
1 Note that the double dots denote normal ordering. See Definition 4.4<br />
2 We write it like this to show what it would be if xn = 0.<br />
49<br />
j=1<br />
Θ(x 0 j) 2 .
In our setting this is not defined, but it might help us get an idea of what<br />
to do.<br />
The discontinuity plane of χn(x) is<br />
Pχn = {x = (x1, . . . , xn−1, 0)|x 0 j = 0 for all j}.<br />
Now if y ∈ supp d ⊂ Γ + n (0) ∪ Γ− n (0), then y2 j<br />
then y0 j = 0 for all j hence<br />
−<br />
3<br />
(y i j) 2 ≥ 0<br />
i=1<br />
≥ 0 for all j. If also y ∈ Pχn,<br />
This is only possible if yi = 0 for i = 1, 2, 3. Hence the intersection is the<br />
origin, that is, supp d ∩ Pξn = {0}.<br />
We now define some tools to investigate the singularity at the origin.<br />
Definition 7.1. The distribution d(x) ∈ S ′ (R m ) is said to have a<br />
quasi-asymptotics d0(x) at x = 0 with respect to a positive continuous<br />
function ρ(δ), δ > 0, if the limit<br />
exists in S ′ (R m ) 3 .<br />
lim<br />
δ→0 ρ(δ)δmd(δx) = d0(x) = 0 (7.1)<br />
The equivalent definition in momentum space reads<br />
Definition 7.2. The distribution ˆ d(p) ∈ S ′ (R m ) has quasi-asymptotics<br />
ˆd0(p) at p = ∞ if<br />
exists for all ˇ φ ∈ S(R m ).<br />
lim<br />
δ→0 ρ(δ)〈 ˆ d( p<br />
δ ), ˇ φ(p)〉 = 〈 ˆ d0, ˇ φ〉 = 〈0, ˇ φ〉 (7.2)<br />
We should show these definitions are indeed equivalent. This follows since<br />
δ m 〈d(δx), φ(x)〉 = 〈d(x), λδφ(x)〉<br />
= 〈 ˆ d(p), (λδφ)ˇ(p)〉<br />
= 〈 ˆ d(p), δ m (λ 1/δ ˇ φ)(p)〉, by (7.4)<br />
= 〈 ˆ d( p<br />
δ ), ˇ φ(p)〉, (7.3)<br />
by Proposition 1.12 and using<br />
(λδφ)ˇ(p) = (2π) −m<br />
<br />
e ix·p <br />
x<br />
<br />
φ dx = (2π)<br />
δ<br />
−m δ m<br />
<br />
e iy·δp φ(y)dy<br />
= δ m λ 1/δ ˇ φ(p). (7.4)<br />
3 An explanation of why we require d0 = 0 is found after Definition 7.4<br />
50
Let’s show that quasi-asymptotics do in fact say something about the<br />
origin, only.<br />
Say, d(x) = d1(x) + d2(x) where supp d1 is compact containing {0} and<br />
supp d2 is bounded away from {0}. Then<br />
lim<br />
δ→0 ρ(δ)δm 〈d2(δx), φ0(x)〉 = lim ρ(δ)〈d2(x), φ0(<br />
δ→0 x<br />
)〉 = 0,<br />
δ<br />
for all φ0 ∈ C ∞ 0 (Rm ) due to Proposition 1.12 and the support properties of<br />
d2. Since C ∞ 0 is dense in S the distribution also vanishes on S, hence<br />
lim<br />
δ→0 ρ(δ)δm 〈d1(δx), φ(x)〉 = 〈d0, φ〉 for all φ ∈ S.<br />
FIXME: why does K0 need to be compact?<br />
Lemma 7.3. Let ρ be the positive continuous function of definitions 7.1<br />
and 7.2. Then there exists ω ∈ R such that<br />
for all a > 0.<br />
Proof. By Proposition 1.12<br />
ρ(aδ)<br />
lim<br />
δ→0 ρ(δ) = aω , (7.5)<br />
lim<br />
δ→0 ρ(δ)〈 ˆ d( p<br />
δ ), (λ1/a ˇ φ)(p)〉 = a −m lim ρ(δ)〈<br />
δ→0 ˆ d( p<br />
aδ ), ˇ φ(p)〉<br />
= a −m lim<br />
δ→0<br />
= a −m lim<br />
δ→0<br />
ρ(δ)<br />
ρ(aδ) ρ(aδ)〈 ˆ d( p<br />
ρ(δ)<br />
ρ(aδ) 〈 ˆ d0(p), ˇ φ(p)〉,<br />
since the limit<br />
lim<br />
δ→0 ρ(aδ)〈 ˆ d( p<br />
aδ ), ˇ φ(p)〉 = 〈 ˆ d0(p), ˇ φ(p)〉<br />
exists. Thus the following is defined<br />
aδ ), ˇ φ(p)〉<br />
ρ(aδ)<br />
ρ0(a) := lim<br />
δ→0 ρ(δ) = a−m 〈 ˆ d0(p), ˇ φ(p)〉<br />
〈 ˆ d0(p), ˇ . (7.6)<br />
φ(ap)〉<br />
Note that ρ0 is continuous, since λ 1/a : S → S is continuous, and that<br />
ρ(aδ)<br />
ρ0(a)ρ0(b) = lim<br />
δ→0 ρ(δ) lim<br />
ρ(bδ) ρ(abδ)<br />
= lim<br />
δ→0 ρ(δ) δ→0 ρ(bδ) lim<br />
ρ(bδ) ρ(abδ)<br />
= lim<br />
δ→0 ρ(δ) δ→0 ρ(δ)<br />
Now define a function by f(s) = ln(ρ0(e s )). Then<br />
= ρ0(ab).<br />
f(s) + f(t) = ln(ρ0(e s )) + ln(ρ0(e t )) = ln(ρ0(e t )ρ0(e s )) = ln(ρ0(e s+t ))<br />
51<br />
= f(s + t).
Let n, m ∈ N then<br />
and<br />
Hence<br />
mf( n n<br />
n<br />
) = f( ) + · · · + f(<br />
m m m ) = f(m<br />
<br />
m terms<br />
n<br />
) = f(n)<br />
m<br />
f(n) = nf(1).<br />
f( n n<br />
) =<br />
m m f(1).<br />
Since ρ0 is continuous f(x) = xf(1) for all x > 0.<br />
Thus<br />
ρ0(e x ) = e f(x) = e xf(1) ,<br />
that is, substituting x = ln a,<br />
Let ω := f(1).<br />
ρ0(a) = e f(1) ln(a) ln af(1)<br />
= e = a f(1) .<br />
Note that (7.5) is the condition of a regularly varying function.<br />
We regard ω as an indication of how singular a causal distribution is near<br />
{0}. This leads us to define:<br />
Definition 7.4. The distribution d ∈ S ′ (R m ) is called singular of order ω,<br />
if it has a quasi-asymptotics d0(x) at x = 0, or its Fourier transform has<br />
quasi-asymptotics ˆ d0(p) at p = ∞, respectively, with power-counting<br />
function ρ(δ) satisfying<br />
for each a > 0.<br />
ρ(aδ)<br />
lim<br />
δ→0 ρ(δ) = aω ,<br />
Note that d0 = 0 in (7.1) and the corresponding requirement in (7.2),<br />
respectively, are required to make sure that the singular order is uniquely<br />
defined. If we skipped the requirement any ω ′ ≥ ω could be chosen as<br />
singular order.<br />
Fixme: It this explained well enough?<br />
Lemma 7.5. The distributions d0 and ˆ d0 are homogeneous of degree<br />
−(m + ω) and ω respectively.<br />
Proof. From (7.6) and (7.5),<br />
Hence by Proposition 1.12<br />
a −ω 〈 ˆ d0(p), ˇ φ(p)〉 = a m 〈 ˆ d0(p), ˇ φ(ap)〉.<br />
〈 ˆ d0( p<br />
a ), ˇ φ(p)〉 = a m 〈 ˆ d0(p), ˇ φ(ap)〉 = a −ω 〈 ˆ d0(p), ˇ φ(p)〉. (7.7)<br />
52
Thus the condition of Definition 1.14 is satisfied by ˆ d0.<br />
Further from (7.3)<br />
and<br />
hence by (7.7)<br />
a m 〈d0(ax), φ(x)〉 = 〈d0(x), φ( x<br />
a )〉 = 〈 ˆ d0( p<br />
a ), ˇ φ(p)〉.<br />
〈 ˆ d0(p), ˇ φ(p)〉 = 〈d0(x), φ(p)〉<br />
〈d0(ax), φ(x)〉 = a −(m+ω) 〈d0(x), φ(x)〉.<br />
The condition of Definition 1.14 is satisfied by d0.<br />
The singular order now gives us a way to distinguish two different<br />
situations in the splitting process. We investigate the splitting first in the<br />
case where ω < 0 and then for ω ≥ 0.<br />
53
7.2 Case I: Negative Singular Order<br />
7.2.1 Existence<br />
In this subsection we will show that there exists a decomposition of d into<br />
an advanced part a and a retarded part r.<br />
Choose ɛ > 0 such that ω + ɛ < 0 then t ω+ɛ → ∞ for t → 0 − hence by<br />
(6.11) ρ(t) → ∞ for t → 0. Hence by (7.1)<br />
lim<br />
δ→0<br />
〈d(x), φ(x<br />
δ<br />
)〉 = lim<br />
δ→0 δm 〈d0(x), φ(x)〉<br />
〈d(δx), φ(x)〉 = lim<br />
= 0. (7.8)<br />
δ→0 ρ(δ)<br />
Let v = (v1, . . . , vn−1) ∈ Γ + n−1 (0) be an arbitrary but fixed vector and<br />
define a hyperplane by<br />
PH := {x ∈ M n−1 n−1 <br />
|v · x := 〈vj, xj〉 = 0}.<br />
Lemma 7.6. All products satisfy<br />
and<br />
j=1<br />
〈vj, xj〉 ≥ 0 if x ∈ Γ + n−1 (0) (7.9)<br />
〈vj, xj〉 ≤ 0 if x ∈ Γ − n−1<br />
(0). (7.10)<br />
Proof. For all j, (v0 j )2 − 3 i=1 (vi j )2 ≥ 0 hence v0 j ≥<br />
<br />
3<br />
i=1 (vi j )2 since<br />
v0 j ≥ 0. Now if x ∈ Γ+ n−1 (0), then for all j<br />
〈vj, xj〉 = v 0 j x 0 j −<br />
3<br />
i=1<br />
v i jx i j ≥ v 0 j x 0 j −<br />
<br />
<br />
<br />
3 <br />
i=1<br />
If on the other hand x ∈ Γ − n−1 (0), then x0 j<br />
〈vj, xj〉 = v 0 j x 0 j −<br />
3<br />
i=1<br />
(v i j )2<br />
v i jx i j ≤ v 0 j x 0 j − x 0 j<br />
≤ x 0 j<br />
FIXME: read this through carefully.<br />
<br />
<br />
<br />
3 <br />
Thus PH splits causal support (5.20) of d.<br />
54<br />
i=1<br />
3<br />
i=1<br />
(x i j )2 ≥ v 0 j x 0 j − v 0 j x 0 j = 0<br />
≤ 0 for all j, hence<br />
<br />
<br />
<br />
3 <br />
i=0<br />
(v i j )2 − x 0 j<br />
(v i j )2<br />
<br />
<br />
<br />
3 <br />
i=0<br />
(v i j )2 = 0
Now choose a monotonous function χ0 ∈ C∞ (R) such that<br />
⎧<br />
⎪⎨ 0, t ≤ 0,<br />
χ0(t) = s ∈ [0, 1),<br />
⎪⎩<br />
1,<br />
0 < t < 1,<br />
t ≥ 1.<br />
The following theorem gives us a retarded distribution.<br />
Theorem 7.7. The limit<br />
exists.<br />
v · x def<br />
lim χ0( )d(x) = θ(v · x)d(x) (7.11)<br />
δ→0 δ<br />
In order to prove this we need the following lemma.<br />
Lemma 7.8. Let a1 > 1, then<br />
ψ0( x<br />
)d(x) :=<br />
δ<br />
uniformly in a ≥ a1.<br />
<br />
χ0(a<br />
v · x v · x<br />
) − χ0(<br />
δ δ )<br />
<br />
d(x) → 0 for δ → 0<br />
Proof. Let U be a neighborhood of Γ + (0) ∪ Γ − (0). Choose a function<br />
ψ1 ∈ C ∞ (R m ) such that<br />
Then, by (5.20)<br />
⎧<br />
⎪⎨<br />
ψ1(x) =<br />
⎪⎩<br />
0, Rm \ U.<br />
〈ψ0( x<br />
δ<br />
1, x ∈ Γ + (0) ∪ Γ − (0),<br />
s ∈ (0, 1], U \ Γ + (0) ∪ Γ − (0),<br />
x x<br />
)d(x), φ(x)〉 = 〈ψ0( )d(x), ψ1(<br />
δ δ )φ(x)〉<br />
= 〈φ(x)d(x), ψ0( x<br />
δ<br />
for any φ ∈ S(R m ), since ψ0ψ1 ∈ C ∞ 0 (Rm ).<br />
By (7.1)<br />
dδ def<br />
= ρ(δ)δ m d(δx) → d0<br />
for δ → 0. Further<br />
in C ∞ (R m ) since<br />
φ(δx) → φ(0) for δ → 0<br />
D α (φ(δx) − φ(0)) = δ |α| (D α φ)(δx) − D α φ(0).<br />
Therefore by [12] Theorem 6.1.8<br />
x<br />
)ψ1( )〉, (7.12)<br />
δ<br />
φ(δx)dδ → φ(0)d0 for δ → 0. (7.13)<br />
55
Therefore φ(x)d(x) also has quasi-asymptotics of order ω < 0 at x = 0<br />
with respect to ρ.<br />
Hence by (7.8) and (7.12)<br />
〈ψ0( x<br />
)d(x), φ(x)〉 → 0 for δ → 0 (7.14)<br />
δ<br />
for all φ ∈ S(R m ).<br />
Now for fixed a1 > 1 the convergence is uniform on the compact interval<br />
1 ≤ a ≤ a1.<br />
We want to show the uniformity of the limit in a ′ ≥ a1 > 1.<br />
To this end, we claim that<br />
Claim.<br />
χ0<br />
<br />
n v · x<br />
<br />
v · x<br />
n−1 <br />
a − χ0 =<br />
δ<br />
δ<br />
j=0<br />
ψ0<br />
<br />
j x<br />
<br />
a<br />
δ<br />
We prove this by induction. It holds for n = 1 by definition. Assume it<br />
holds for n = k − 1. Then,<br />
<br />
k v · x<br />
<br />
v · x<br />
<br />
χ0 a − χ0<br />
<br />
δ<br />
δ<br />
k v · x<br />
<br />
k−1 v · x<br />
<br />
k−1 v · x<br />
<br />
v · x<br />
<br />
= χ0 a − χ0 a + χ0 a − χ0<br />
δ<br />
δ<br />
δ<br />
δ<br />
<br />
k−1 x<br />
k−2<br />
<br />
j x<br />
<br />
= ψ0 a + ψ0 a .<br />
δ<br />
δ<br />
j=0<br />
Which proves the claim.<br />
We write a ′ = a n , for a suitable n and 1 ≤ a ≤ a1. Now the problem is<br />
translated into showing uniformity in n. By (7.13) we may choose δ0 such<br />
that for δ < δ0<br />
|ρ(δ)〈φ(x)d(x), ψ1(x/δ)ψ0(x/δ)〉| ≤ |〈φ(0)d0(x), ψ1(x)ψ0(x)〉| + 1 = K.<br />
Note that δ/a j ≤ δ < δ0. Thus<br />
From (7.5) we know that<br />
Now since a −ω > a −ω<br />
1<br />
Since δ/a q < δ for all q ≥ 1<br />
|〈φ(x)d(x), ψ1(a j x/δ)ψ0(a j x/δ)〉| ≤ K/ρ(δ/a j )<br />
ρ(δ/a)<br />
ρ(δ) .<br />
for ω ≤ 0 we may assume, that for δ < δ0<br />
ρ(δ/a) ≥ a −ω<br />
1 ρ(δ).<br />
ρ(δ/a j ) ≥ a −ω<br />
1 ρ(δ/aj−1 ) ≥ a 2ω<br />
1 ρ(δ/a j−2 ) ≥ a −jω<br />
1<br />
56<br />
ρ(δ).
Thus<br />
n−1 <br />
|<br />
j=0<br />
〈φ(x)d(x), ψ1(a j x/δ)ψ0(a j n−1<br />
x/δ)〉| ≤<br />
But ρ(δ) → ∞ Hence the convergence is uniform.<br />
We now prove the theorem.<br />
<br />
K/ρ(δ/a j )<br />
j=0<br />
≤ Kρ(δ) −1<br />
∞<br />
j=0<br />
a jω<br />
1<br />
= Kρ(δ) −1 (1 − a ω 1 ) −1 .<br />
Proof of Theorem 7.7. Another way to state the theorem is that the<br />
limit in (7.11) exists for all sequences δn → 0. This is done by showing<br />
〈χ0( v·x)d(x),<br />
φ(x)〉 is a Cauchy sequence. Consider<br />
δn<br />
v · x v · x<br />
v · x v · x<br />
χ0( ) − χ0( ) = χ0(an,m ) − χ0( ), (7.15)<br />
δm<br />
δn<br />
where an,m = δn/δm. Assuming δm < δn, we know from Lemma 7.8 that<br />
for all φ and for all ɛ > 0 there exists some δ0 > 0 such that for all δ < δ0<br />
|〈(χ0(an,m<br />
v · x<br />
δ<br />
δn<br />
δn<br />
v · x<br />
) − χ0( ))d(x), φ(x)〉| < ɛ<br />
δ<br />
Hence given φ and ɛ > 0 we may choose δ0 and N such that δn < δ0 for all<br />
n, m ≥ N<br />
v · x v · x<br />
|〈(χ0( ) − χ0( ))d(x), φ(x)〉| < ɛ.<br />
δm<br />
δn<br />
Hence 〈χ0( v·x)d(x),<br />
φ(x)〉 is a Cauchy sequence.<br />
δn<br />
Now suppose δn and δn ′ are two sequences converging towards 0. Then as<br />
above, with an,n ′ = δn ′/δn in (7.15),<br />
v · x<br />
v · x<br />
〈χ0( )d(x), φ(x)〉 − 〈χ0( )d(x), φ(x)〉 → 0<br />
δn<br />
as n → ∞. Hence the limit exists.<br />
Now we may define<br />
Definition 7.9.<br />
δn ′<br />
v · x<br />
r(x) = lim χ0( )d(x) = θ(v · x)d(x).<br />
δ→0 δ<br />
Similarly there exists an advanced distribution.<br />
57
Corollary 7.10. The limit<br />
exists.<br />
−v · x<br />
lim χ0( )d(x) := −a(x)<br />
δ→0 δ<br />
Proof. This follows by the proof of Lemma 7.8 and Theorem 7.7 with χ0(t)<br />
substituted by χ0(−t).<br />
FIXME: think about this.<br />
Note that by (7.9) and (7.10),<br />
and<br />
Hence<br />
v · x<br />
supp r ⊂ supp(lim χ0(<br />
δ→0 δ )) ⊂ Γ+ n−1 (0)<br />
supp a ⊂ supp(lim<br />
δ→0 χ0(<br />
−v · x<br />
)) ⊂ Γ<br />
δ<br />
− n−1 (0).<br />
r(x) − a(x) = Θ(v · x)d(x) − Θ(−v · x)d(x) = d(x). (7.16)<br />
These support properties allows us to apply a and r to discontinuous<br />
test-functions as follows<br />
〈r, φ〉 = 〈r(x), Θ(v · x)φ(x)〉 and 〈r, (1 − Θ)φ〉 = 0<br />
〈a, φ〉 = 〈a(x), (1 − Θ)(v · x)φ(x)〉 and 〈a, Θφ〉 = 0.<br />
Thus by (7.16) we may split d in a retarded and an advanced part, simply<br />
by multiplication with Θ,<br />
and<br />
Note that by (7.1)<br />
and<br />
〈d, Θφ〉 = 〈r, Θφ〉 − 〈a, Θφ〉 = 〈r, φ〉<br />
〈d, (1 − θ)φ〉 = 〈r, (1 − Θ)φ〉 − 〈a, (1 − Θ)φ〉 = −〈a, φ〉.<br />
〈d0, Θφ〉 = lim<br />
δ→0 ρ(δ)δ m 〈d(δx), Θφ(x)〉 = lim<br />
δ→0 ρ(δ)〈d(x), Θλδφ(x)〉<br />
= lim ρ(δ)〈r(x), φ(x/δ)〉<br />
δ→0<br />
= lim ρ(δ)δ<br />
δ→0 m 〈r(δx), φ(x)〉<br />
〈d0, (1 − Θ)φ〉 = lim<br />
δ→0 ρ(δ)δ m 〈d(δx), (1 − Θ)φ(x)〉 = lim<br />
δ→0 ρ(δ)〈a(x), φ(x/δ)〉<br />
We see that r and a have same singular order as d.<br />
58<br />
= lim<br />
δ→0 ρ(δ)δ m 〈a(δx), φ(x)〉
7.2.2 Uniqueness<br />
We still need to prove that the decomposition of d is unique. To this end<br />
we need the following result<br />
Theorem 7.11. Let u ∈ D ′ (R n ) with support supp u = {0}. Then there<br />
exists a N ∈ N such that<br />
where cα ∈ C.<br />
u = <br />
|a|≤N<br />
cα∂ α δ,<br />
This theorem is a consequence of the following lemmas.<br />
Lemma 7.12.<br />
1<br />
k! ∂k t f(xt) = <br />
|α|=k<br />
x α<br />
α! f (α) (xt) (7.17)<br />
Proof. The first we note that by the chain rule, we may write the left hand<br />
side as<br />
1<br />
k! ∂k t f(xt) = 1<br />
<br />
n<br />
k <br />
∂ <br />
xi f(y) <br />
k! ∂yi<br />
i=1<br />
y=xt<br />
It is this form we will use when proving (7.17) by induction. That is, we<br />
want to prove<br />
1<br />
n<br />
∂<br />
k xi f(y) =<br />
k! ∂yi<br />
xα α! ∂αf(y) i=1<br />
59<br />
|α|=k
It is obvious for k = 0 and k = 1. Now, assume it holds for k = m. Then<br />
1<br />
n<br />
(m + 1)!<br />
= 1<br />
n<br />
m + 1<br />
j=1<br />
= 1<br />
n<br />
m + 1<br />
= 1<br />
m + 1<br />
= 1<br />
m + 1<br />
= 1<br />
m + 1<br />
= 1<br />
m + 1<br />
= <br />
j=1<br />
n<br />
i=1<br />
xj<br />
xj<br />
<br />
xi<br />
∂<br />
∂yj<br />
∂<br />
∂yj<br />
∂<br />
∂yi<br />
m+1<br />
f(y)<br />
<br />
1<br />
n<br />
m!<br />
<br />
x<br />
xj<br />
j=1 |α|=m<br />
α<br />
α!<br />
n<br />
<br />
j=1 |α|=m<br />
<br />
|α|=m j=1<br />
<br />
|β|=m+1 j=1<br />
n<br />
|β|=m+1 j=1<br />
|α|=m<br />
i=1<br />
xi<br />
∂<br />
∂yi<br />
x α<br />
α! ∂α f(y)<br />
∂<br />
∂<br />
∂yj<br />
α f(y)<br />
m<br />
f(y)<br />
xα+(0,...,0,1,0,...,0) ∂<br />
α!<br />
α+(0,...,0,1,0,...,0) f(y)<br />
n<br />
x<br />
(αj + 1)<br />
α+(0,...,0,1,0,...,0)<br />
(α + (0, . . . , 0, 1, 0, . . . , 0))! ∂α+(0,...,0,1,0,...,0) f(y)<br />
n<br />
x β<br />
β! ∂β f(y)<br />
x<br />
βj<br />
β<br />
β! ∂βf(y), where β = α + (0, ..., 0, 1, 0, ...0)<br />
where only the jth entry of (0, . . . , 0, 1, 0, . . . , 0) is non-zero and where we<br />
have used<br />
(α + (0, . . . , 0, 1, 0, . . . , 0))! = (α1, α2, . . . , αj + 1, . . . , αn)!<br />
Finally, rename β = α.<br />
= α1!α2! · · · (αj + 1)! · · · αn!<br />
= (αj + 1)α!.<br />
Lemma 7.13. Let f ∈ C ∞ 0 (Rn ) and t ∈ R, then<br />
f(x) ≤ <br />
for any N ≥ 1.<br />
|α|≤N−1<br />
xα α! f (α) <br />
N N<br />
(0) + |x|<br />
α!<br />
|α|=N<br />
sup |f (α) (x ′ )| <br />
′<br />
|x | < |x| ,<br />
60
Proof. Let F : t ↦→ f(xt). By application of Taylor’s formula [10] to F ,<br />
f(x) = F (1) =<br />
=<br />
N−1 <br />
k=0<br />
N−1 <br />
k=0<br />
= <br />
1<br />
k! ∂k t F (0) +<br />
<br />
|α|≤N−1<br />
|α|=k<br />
1<br />
(N − 1)!<br />
xα α! f (α) <br />
(0) + N<br />
xα α! f (α) (0) + <br />
|α|=N<br />
1<br />
(1 − t)<br />
0<br />
N−1 ∂ N t F (t)dt<br />
1<br />
N−1<br />
(1 − t)<br />
xα α! f (α) (xt)dt<br />
0<br />
|α|=N<br />
x α<br />
α! fα(x), (7.18)<br />
where fα(x) = N 1<br />
0 (1 − t)N−1 f (α) (xt)dt and where we have used (7.17).<br />
Now, since |x α | ≤ maxj |xj| N ≤ |x| N when |α| = N, we can make the<br />
following estimate<br />
As wanted.<br />
<br />
|α|=N<br />
xα α! fα(x)<br />
<br />
N 1<br />
≤ |x|<br />
α!<br />
|α|=N<br />
fα(x)<br />
<br />
N<br />
= |x|<br />
|α|=N<br />
<br />
N<br />
≤ |x|<br />
|α|=N<br />
<br />
N<br />
≤ |x|<br />
|α|=N<br />
1<br />
α! N<br />
1<br />
(1 − t)<br />
0<br />
N−1 f (α) (xt)dt<br />
N<br />
α! sup |(1 − t)<br />
t∈[0,1]<br />
N−1 f (α) (xt)|<br />
N<br />
α! sup |f (α) (x ′ )| |x ′ | < |x| .<br />
Lemma 7.14. Let u ∈ D(R n ) with support supp u = 0, then there exists<br />
an N ≥ 0 such that 〈u, φ〉 = 0 for all φ ∈ C ∞ 0 (Rn ) and ∂ α φ(0) = 0 when<br />
|α| ≤ N.<br />
Proof. Let ψ ∈ C ∞ 0 (Rn ) be defined such that<br />
Let ɛ ∈ (0, 1). Then<br />
⎧<br />
⎪⎨ 1, |x| < 1/2,<br />
ψ(x) = s ∈ [0, 1],<br />
⎪⎩<br />
0,<br />
1/2 ≤ |x| ≤ 1,<br />
|x| > 1.<br />
〈u, φ〉 = 〈u(x), φ(x)ψ(x/ɛ)〉, for all φ ∈ C ∞ 0 (R n ),<br />
since supp u ⊂ {x ∈ R n ||x| < ɛ/2} =: U and ψ(x/ɛ) = 1 on U.<br />
If x > 1, then x/ɛ > x > 1, hence the map x ↦→ φ(x)ψ(x/ɛ) is supported in<br />
the unit-ball. From now on we may therefore assume that |x| ≤ ɛ.<br />
61
By compactness of the unit-ball and the definition of distributions there<br />
exists a real number C ≥ 0 and a non-negative integer N such that<br />
|〈u, φ〉 ≤ C <br />
α≤N<br />
sup |∂ α (φ(x)ψ(x/ɛ))|, for all φ ∈ C ∞ 0 (R n ). (7.19)<br />
If ∂ α φ(0) = 0 when |α| ≤ N and |β| ≤ N, we may apply Lemma 7.13 to<br />
∂ β φ as follows<br />
|∂ β φ(x)| ≤ |x| N−β+1<br />
≤ ɛ N+1−β<br />
By Leibniz’ formula,<br />
<br />
|γ|=N+1−β<br />
<br />
|γ|=N+1−β<br />
∂ α (φ(x)ψ(x/ɛ))) = <br />
α=β+γ<br />
= <br />
α=β+γ<br />
N + 1 − β<br />
γ!<br />
N + 1 − β<br />
γ!<br />
sup |∂ γ (x ′ )| |x ′ | < |x| <br />
sup |∂ γ (x)| |x| < 1 . (7.20)<br />
α!<br />
β!γ! ∂β φ(x)∂ γ (ψ(x/ɛ))<br />
α!<br />
β!γ! ɛ−|γ| ∂ β φ(x)(∂ γ ψ)(x/ɛ). (7.21)<br />
Combining (7.20) and (7.20) together with the fact that ∂ γ ψ(x/ɛ has<br />
compact supported contained in the unit-ball, we see that there exists<br />
constants Cα for each α ≤ N which is independent of ɛ such that<br />
FIXME: (start) Could say that supp ∂ γ compact (in unit-ball) so attains<br />
its maximum value. For a fixed α there is only finitely many possible γ, so<br />
the maximum of these supremums gives us the estimate. Fixme: (end)<br />
Plugging this into (7.19),<br />
|∂ α (φ(x)ψ(x/ɛ)) ≤ Cαɛ N+1−|β| e −|γ| = Cαɛ N+1−|α|<br />
|〈u, φ〉| ≤ C <br />
|α|≤N<br />
letting K be the constant. Finally, let ɛ → 0.<br />
FIXME: maybe this could be better.<br />
Cαɛ N+1−|α| ≤ Kɛ,<br />
Proof of Theorem 7.11. Pick N and ψ as in Lemma 7.14.<br />
By (7.18) every φ ∈ C ∞ 0 (Rn ) can be written on the form<br />
φ = ψ(x) <br />
|α|≤N<br />
62<br />
x α<br />
α! ∂α φ(0) + φ ′
with N and ψ as in the lemma. Then φ ′ ∈ C ∞ 0 (Rn ) and ∂ α φ ′ (0) = 0 for<br />
|α| ≤ N. Hence, by Lemma 7.14 〈u, φ ′ 〉 = 0. Thus<br />
〈u, φ〉 = 〈u(x), ψ(x) <br />
|α|≤N<br />
= 〈u(x), ψ(x) <br />
= <br />
|α|≤N<br />
= <br />
|α|≤N<br />
with cα = (−1) |α| 〈u(x),ψ(x)xα 〉<br />
α! .<br />
|α|≤N<br />
x α<br />
α! ∂α φ(0) + φ ′ 〉<br />
x α<br />
α! ∂α φ(0)〉<br />
〈u(x), ψ(x)xα 〉<br />
∂<br />
α!<br />
α φ(0)<br />
cα∂ α δ,<br />
Now to prove the decomposition d = r − a is unique assume there are two<br />
solutions {r1, a1} and {r2, a2} to the splitting problem. That is,<br />
d = r1 − a1 = r2 − a2 hence r1 − r2 = a1 − a2.<br />
The support of the differences is {0}<br />
FIXME: Is it??<br />
hence by Theorem 7.11<br />
r1 − r2 = a1 − a2 = <br />
cα∂ α δ. (7.22)<br />
|α|≤N<br />
We want to investigate each term of the sum. From (8.6) (FIXME: maybe<br />
better to place this equation before) we know that ˆ δ = 1, hence by<br />
Theorem 5.17 in [2]<br />
Hence,<br />
ρ(δ)〈 cα∂ α δ( p<br />
δ ), ˇ φ(p)〉 = ρ(δ)cαi |α| 〈<br />
∂ α δ = i |α| D α δ = i |α| p αˆ δ = i |α| p α<br />
<br />
p<br />
α ,<br />
δ<br />
ˇ φ(p)〉 = ρ(δ)<br />
δ |α| 〈 cα∂ αδ(p), ˇ φ(p)〉.<br />
Thus ρ(δ) = δ |α| for dα(x) = cα∂ α δ(x) by (7.2) and then clearly the<br />
singular order is ωα = |α|.<br />
For r1 − r2 = <br />
|α|≤N cα∂ α δ we then must have |α| = ωα ≤ ω for all α.<br />
This implies<br />
r1 − r2 = <br />
|α|≤ω<br />
cα∂ α δ (7.23)<br />
Further in our case ω < 0. Hence all the cα must vanish in (7.22). This<br />
shows r1 = r2 and a1 = a2. Hence the splitting is unique.<br />
63
7.3 Case II: Positive Singular Order<br />
Choose positive ɛ < 1 then since C ′ δ ɛ−1 = C′ δ ω+ɛ<br />
δ ω+1<br />
≤ ρ(δ)<br />
δ ω+1 by (6.11),<br />
ρ(δ)<br />
→ ∞ for δ → 0<br />
δω−1 Hence if we choose a multi-index b such that |b| = ω + 1, then<br />
lim<br />
δ→0 〈d(x)xb , ψ(x/δ)〉 = lim δ<br />
δ→0 m+ω+1 〈d(δx)x b , ψ(x)〉<br />
= lim<br />
δ→0<br />
δ m+ω<br />
ρ(δ) 〈d0(x), x b ψ(x)〉 = 0, (7.24)<br />
by (7.1).<br />
We want to show the splitting can be done for test-functions φ satisfying<br />
Definition 7.15. We define 4<br />
D a φ(0) = 0 for |a| ≤ ω.<br />
S ω (R m ) = {φ ∈ S(R m )|D α φ(0) = 0 for all α with |α| ≤ ω}<br />
This is the subspace of S(R m ) of test-functions which vanish up to order ω<br />
at 0.<br />
FIXME: show ∞ > sing order(d) = ω ≥ 0 ⇒ d ∈ S ′ω (R m ).<br />
Definition 7.16. We define operators W (ω,w) : S(R n ) → S ω (R n ) by<br />
(W (ω,w)φ)(x) := φ(x) − w(x)<br />
for each w(x) ∈ S(R n ) such that<br />
ω<br />
|a|=0<br />
w(0) = 1 and D α w(0) = 0 for 1 ≤ |a| ≤ w.<br />
x a<br />
a! (Da φ)(0) (7.25)<br />
Recognizing Taylor’s formula given by (7.18) we note that W is actually a<br />
modified Taylor subtraction operator which projects S(Rm ) into the space<br />
of test-functions which vanish up to order ω at 0.<br />
Note that we may write<br />
(W φ)(x) = <br />
x b ψb(x)<br />
|b|=ω+1<br />
by (7.18)<br />
We now want to apply θ from Theorem 7.7.<br />
4 If ω is not an integer use [ω] the largest integer such that [ω] < ω. The reason I<br />
haven’t written it like this straight away, is that ω always seems to be an integer in QED.<br />
64
Lemma 7.17. Let d (b) (x) = x b d(x), then d (b) has quasi-asymptotics<br />
x b d0(x) with respect to ρ (b) (δ) = δ −|b| ρ(δ), and it has singular order ω − |b|.<br />
Proof. The quasi-asymptotics follows from<br />
lim<br />
δ→0 ρ(b) δ m d (b) (δx) = lim<br />
delta→0 ρ(b) (δ)δ m+|b| x b d(δx) = lim<br />
δ→0 ρ(δ)δ m x b d(δx)<br />
The singular order is given by<br />
for all a > 0.<br />
= x b d0(x)<br />
ρ<br />
lim<br />
δ→0<br />
(b) (aδ)<br />
ρ (b) (δa)<br />
= lim<br />
(δ) δ→0<br />
−|b| ρ(aδ)<br />
δ−|b| = aω−|b|<br />
ρ(δ)<br />
By the previous lemma x b d(x) has singular order ω − |b|. Hence it is<br />
negative for b ≥ ω + 1. Thus we may define<br />
Definition 7.18.<br />
and<br />
〈r(x), φ〉 := 〈d, θ(v · x)W φ〉<br />
a(x) := d(x) − r(x).<br />
By construction supp r ⊂ Γ + n−1<br />
x ∈ Γ + n−1 (0) \ {0} by (7.25), since a test-function φ ∈ S supported in<br />
Γ + n−1 (0) \ {0} satisfies Dαφ(0) = 0 for all α.<br />
Like for the case ω < 0, r and a have the same singular order as d. This<br />
follows from the equations<br />
(0). Further r(x) = d(x) for<br />
lim ρ(δ)〈r(x), φ(x/δ)〉 = lim ρ(δ)〈d(x), (θW φ)(x/δ)〉<br />
δ→0 δ→0<br />
= 〈d0(x), θW φ(x)〉.<br />
In contrast to the case where ω < 0 the splitting is not unique for w ≥ 0<br />
because of the dependence on w. Again the support of the difference<br />
between two solutions is {0}, hence by (7.23)<br />
r1 − r2 = <br />
Cα∂ α δ<br />
|α|≤ω<br />
The coefficients Cα have to be fixed by additional normalization conditions.<br />
65
Chapter 8<br />
Application to QED<br />
8.1 Using the Game Plan<br />
In this section we will use the game plan from the end of Section 5.1 to<br />
find D2 from T1 and then show how to find T2.<br />
In QED 1<br />
T1(x) = ie : ψ(x)γ µ ψ(x) : Aµ(x),<br />
By (5.1) this means that<br />
˜T1(x) = −T1(x) = −ie : ψ(x)γ µ ψ(x) : Aµ(x).<br />
We now want to construct an advanced distribution according to (5.7).<br />
This is done using Theorem 4.6 remembering that the only contractions<br />
existing in QED are given by equations (4.1), (4.2) and (4.3).<br />
1 [6] page 183<br />
66
A ′ 2(x1, x2) = ˜ T1(x1)T1(x2) = −T1(x1)T1(x2)<br />
= e 2 γ µ<br />
abγν cd : ψa(x1)ψb(x1) :: ψc(x2)ψd(x2) : Aµ(x1)Aν(x2)<br />
<br />
: ψa(x1)ψb(x1)ψ c(x2)ψd(x2) :<br />
= e 2 γ µ<br />
ab γν cd<br />
= e 2 γ µ<br />
ab γν cd<br />
+ : ψ a(x1) | ψb(x1)ψ c(x2) | ψd(x2) :<br />
+ : | ψ a(x1)ψb(x1)ψ c(x2)ψd(x2) | :<br />
+ : | ψ a(x1) | ψb(x1)ψ c(x2) | ψd(x2) | :<br />
<br />
×<br />
<br />
: ψa(x1)ψb(x1)ψ c(x2)ψd(x2) :<br />
: Aµ(x1)Aν(x2) : + : | Aµ(x1)Aν(x2) | :<br />
+ 1<br />
i S(+)<br />
bc (x1 − x2) : ψa(x1)ψd(x2) :<br />
+ 1<br />
i S(−)<br />
da (x2 − x1) : ψb(x1)ψc(x2) :<br />
− S (+)<br />
bc (x1 − x2)S (−)<br />
da (x2<br />
<br />
− x1)<br />
<br />
×<br />
: Aµ(x1)Aν(x2) : +gµνiD (+)<br />
0 (x1 − x2)<br />
<br />
<br />
<br />
. (8.1)<br />
From (5.8) we see that the retarded distribution is given by (8.1) simply by<br />
substituting x1 ↔ x2. For convenience we also interchange the indices<br />
µ ↔ ν, a ↔ c and b ↔ d, arriving at<br />
R ′ 2(x1, x2) = T1(x2) ˜ T1(x1) = −T1(x2)T1(x1)<br />
<br />
: ψa(x1)ψb(x1)ψ c(x2)ψd(x2) :<br />
= e 2 γ µ<br />
ab γν cd<br />
− 1<br />
i S(+)<br />
da (x2 − x1) : ψb(x1)ψc(x2) :<br />
− 1<br />
i S(−)<br />
bc (x1 − x2) : ψa(x1)ψd(x2) :<br />
− S (+)<br />
da (x2 − x1)S (−)<br />
bc (x1<br />
<br />
− x2)<br />
<br />
×<br />
: Aµ(x1)Aν(x2) : +gµνiD (+)<br />
0 (x2 − x1)<br />
67<br />
<br />
. (8.2)
Finally we can calculate the difference D(x1, x2) by (5.11).<br />
D(x1, x2) = R ′ n(x1, x2) − A ′ n(x1, x2)<br />
= e 2 γ µ<br />
abγν <br />
cd ψa(x1)ψb(x1)ψ c(x2)ψd(x2) : gµνi(D (+)<br />
0 (x2 − x1) − D (+)<br />
0 (x1 − x2))<br />
− 1<br />
<br />
S<br />
i<br />
(+)<br />
da (x2 − x1) + S (−)<br />
da (x2<br />
<br />
− x1) : ψb(x1)ψc(x2) :: Aµ(x1)Aν(x2) :<br />
<br />
− gµν S (+)<br />
da (x2 − x1)D (+)<br />
0 (x2 − x1) + S (−)<br />
da (x2 − x1)D (+)<br />
<br />
0 (x1 − x2) : ψb(x1)ψc(x2) :<br />
− 1<br />
i (S(−)<br />
bc (x1 − x2) + S (+)<br />
bc (x1 − x2))ψa(x1)ψd(x2) :: Aµ(x1)Aν(x2) :<br />
<br />
− gµν S (−)<br />
bc (x1 − x2)D (+)<br />
0 (x2 − x1) + S (+)<br />
bc (x1 − x2)D (+)<br />
<br />
0 (x1 − x2) : ψa(x1)ψd(x2) :<br />
<br />
+ S (+)<br />
bc (x1 − x2)S (−)<br />
da (x2 − x1) − S (−)<br />
bc (x1 − x2)S (+)<br />
da (x2<br />
<br />
− x1) : Aµ(x1)Aν(x2) :<br />
+ 1<br />
i gµν<br />
<br />
S (−)<br />
bc (x1 − x2)S (+)<br />
da (x2 − x1)D (+)<br />
0 (x2 − x1)<br />
<br />
. (8.3)<br />
− S (+)<br />
bc (x1 − x2)S (−)<br />
da (x2 − x1)D (+)<br />
0 (x1 − x2)<br />
This formula may be considered as a correspondent to the Feynman rules.<br />
Each term can be treated separately and represents different scattering<br />
scenarios, where the field operators create and annihilate particles. Like<br />
the Feynman rules, they may be illustrated by graphs. But notice the huge<br />
difference that while all our terms are well-defined distributions, the<br />
Feynman rules are ill-defined for closed loops, which here are represented<br />
by terms with products of the propagators S and D.<br />
Lets us consider the first term. First we note that<br />
D + 0 (x1 − x2) − D + 0 (x2 − x1) is in fact the Jordan-Pauli distribution for<br />
mass m = 0, where the full Jordan-Pauli distribution is given by 2<br />
D(x) =<br />
sgn x0<br />
2π (δ(x2 ) − Θ(x 2 ) m<br />
2 √ x 2 J1(m √ x 2 )).<br />
Note that for the 1-dimensional δ-distribution,<br />
Therefore<br />
〈δ(δ 2 x 2 ), φ(δ 2 x 2 )〉 = δ −2 〈δ(x 2 ), φ(x 2 )〉.<br />
D0(δx) =<br />
sgn x0<br />
2π δ(δ2x 2 sgn x0 δ(x<br />
) =<br />
2π<br />
2 )<br />
δ2 Hence if we choose ρ(δ) = δ 2−n we get the quasi-asymptotics<br />
2 [6] page 89.<br />
lim<br />
δ→0 ρ(δ)δn sgn x0<br />
D0(δx) =<br />
2π δ(x2 ).<br />
68
and for each a > 0,<br />
ρ(aδ)<br />
ρ(δ)<br />
= a2−n<br />
Thus in R 4 , ω = −2. Since the singular order is negative the splitting may<br />
be done by multiplication by the Θ-function (7.16). Hence the retarded<br />
part of the first term of (8.3) is<br />
R2(x1, x2) = −ie 2 : ψ(x1)γ µ ψ(x1)ψ(x2)γ ν ψ(x2) : ΘD0(x1 − x2)<br />
Further R ′ (x1, x2) is given by (8.2) by inspecting each term. By (5.12)<br />
T2(x1, x2)<br />
= R2(x1, x2) − R ′ (x1, x2)<br />
= −ie 2 : ψ(x1)γ µ ψ(x1)ψ(x2)γ ν ψ(x2) :<br />
<br />
ΘD0(x1 − x2) + D (+)<br />
<br />
0 (x2 − x1) .<br />
It is not the aim of this project to introduce Feynman propagators but the<br />
reader with knowledge of these might notice that the expression in<br />
brackets is actually a Feynman propagator describing the exchange of a<br />
photon between electrons.<br />
Figure 8.1: Photon exhange.<br />
8.2 The Adiabatic Limit<br />
In the introduction of the S-matrix we used test-functions g ∈ S(R 4 ),<br />
so-called switching functions. As the name infers they switch off<br />
long-range interaction to prevent infrared divergences 3 . Of course this is<br />
not a good model and we need to consider the so-called adiabatic limit<br />
g → 1 to take long-range interaction like the Coulomb-potential into<br />
account. We show how we may carry out the adiabatic limit. In practise it<br />
can me done by taking the so-called scaling limit.<br />
Let g0 ∈ S(R 4 ) be a fixed test-function such that g0(0) = 1. Then we let<br />
g(x) := g0(ɛx) and take the scaling limit, that is, we let ɛ → 0.<br />
Calculations are in practise usually done in momentum space. With this in<br />
mind we note that<br />
<br />
ˆg(k) =<br />
g(x)e −ix·k d 4 x =<br />
<br />
g0(ɛx)e −ik·x d 4 x = 1<br />
ɛ<br />
k<br />
ˆg0( ). (8.4)<br />
4 ɛ<br />
3 By infrared divergence we simply mean a divergence due to physical phenomena at<br />
very long distances or because of contributions from objects with very small energy<br />
69
Theorem 8.1. Identify ˆg(k) with the distribution it induces, then<br />
Proof. For all φ ∈ S(R 4 )<br />
lim〈ˆg,<br />
φ〉 = lim<br />
ɛ→0 ɛ→0 〈g, ˆ φ〉 = lim<br />
ɛ→0<br />
lim<br />
ɛ→0 ˆg(k) = (2π)4δ(k). <br />
g(k) ˆ φ(k)d 4 <br />
k = lim<br />
ɛ→0<br />
g0(ɛk) ˆ φ(k)d 4 <br />
k<br />
= ˆφ(k)d 4 k, (8.5)<br />
by the Dominated Convergence Theorem.<br />
Now note that generally<br />
〈 ˆ δ, φ〉 = 〈δ, ˆ φ〉 = ˆ <br />
φ(0) = φ(x)dx = 〈1, φ〉. (8.6)<br />
Therefore ˆ δ = 1. Similarly<br />
〈 ˇ δ, φ〉 = 〈δ, ˇ φ〉 = ˇ φ(0) = (2π) −n<br />
<br />
hence<br />
φ(x)dx = (2π) −n 〈1, φ〉,<br />
〈ˆ1, φ〉 = (2π) −n 〈F( ˇ δ), φ〉 = (2π) −n 〈δ, φ〉. (8.7)<br />
That is, ˆ1 = (2π) nδ. It then follows from (8.5), that<br />
<br />
lim〈ˆg,<br />
φ〉 =<br />
ɛ→0<br />
as wanted.<br />
Now by (8.4)<br />
<br />
= 1<br />
ɛ 4n<br />
= 1<br />
ɛ 4n<br />
ˆφ(k)d 4 k = 〈1, ˆ φ〉 = 〈ˆ1, φ〉 = (2π) 4 〈δ, φ〉,<br />
d 4 k1 · · · d 4 kn ˆ T (k1, . . . , kn)ˆg(k1) · · · ˆg(kn)<br />
<br />
<br />
d 4 k1 · · · d 4 kn ˆ T (k1, . . . , kn)ˆg0(k1/ɛ) · · · ˆg0(kn/ɛ)<br />
d 4 k1 · · · d 4 kn ˆ T (ɛk1, . . . , ɛkn)ˆg0(k1) · · · ˆg0(kn),<br />
by substituting ki by ɛki. Finally, we may calculate in the usual way and<br />
in the end take the limit ɛ → 0 without problems.<br />
70
Chapter 9<br />
The Microlocal Approach -<br />
A Condition on the Wave<br />
Front Set.<br />
In this section I will use the notation T ∗ (X) even though X will be an<br />
open subset of R n (and thus T ∗ (X) = X × R n ). I use this notation to<br />
emphasize that the results can be generalized to manifolds 1 . As the results<br />
are based on [4] where we only considered the case where X was open in<br />
R n it would not make sense to consider manifolds.<br />
The Method of Epstein and Glaser may be generalized to work on<br />
manifolds using microlocal analysis.<br />
Definition 9.1. Let Ψ ∈ H. Then the operator-valued function<br />
Ψ(f) def<br />
= : W (f) : Ψ, where W (f)φ def<br />
= exp(iφ(f) ∗∗ ),<br />
is said to be infinitely often differentiable at f = 0 if for each n,<br />
1. For all h ∈ D(M) and all t ∈ R the map t ↦→ Ψ(th) is n times<br />
norm-differentiable at 0.<br />
2. The derivatives at 0 define symmetrical operator-valued distributions.<br />
With Ψ as above we define a distribution<br />
: φ(x1) · · · φ(xn) : def<br />
=<br />
∂n in : W (f) :<br />
∂f(x1) · · · ∂f(xn)<br />
1 In fact, globally hyperbolic manifolds. See [3] page 7.<br />
71<br />
<br />
<br />
<br />
f=0
in the sense that if h ∈ D(M), then<br />
〈: φ(x1) · · · φ(xn) :, Ψ〉(h n <br />
) = : φ(x1) · · · φ(xn) : Ψ h(x1) · · · h(xn) (9.1)<br />
= i −n<br />
<br />
∂nΨ ∂f(x1) · · · ∂f(xn) h(x1) · · · h(xn)<br />
We want to generalize this.<br />
Definition 9.2. The microlocal domain of smoothness D is defined by<br />
where<br />
D = {Ψ ∈ H|Φ(f) is infinitely often differentiable at f = 0<br />
<br />
∂nΨ and for all n ∈ N, WF<br />
∂f n<br />
<br />
(T ∗ M n )± def<br />
={x, k) ∈ T ∗ M|ki ∈ V − (0)}.<br />
⊂ (T ∗ M n )−},<br />
Now we may use the idea of equation (9.1) to define an operator-valued<br />
distribution on D.<br />
Definition 9.3. We call the operator-valued distribution : φ(x1) · · · φ(xn) :<br />
on D defined by (9.1) a Wick monomial. 2<br />
Now we would like to generalize the Wick monomial to a polynomial. That<br />
is, for arbitrary indices l = (l1, . . . , ln) we would like to make sense to<br />
To this end we need the following definition.<br />
: φ l1 (x1) · · · φ ln (xn) : . (9.2)<br />
Definition 9.4. A partial diagonal ∆l1,...,lj , where l1 + · · · + lj = n is a<br />
subset of M n on the form<br />
∆l1,...,lj (M) = {(x1, . . . , x1),<br />
. . . , (xj, . . . , xj)|xi<br />
∈ M, i = 1 . . . , j} M<br />
<br />
l1 times<br />
lj times<br />
j .<br />
Now we want to define the polynomial (9.2) as a restriction of ∂ n Ψ<br />
∂f1···∂fn to<br />
arbitrary partial diagonals. This is done in practice by a pullback. 3<br />
Define the partial diagonal map<br />
by<br />
δn,l : M n → ∆l1,...,ln (M)<br />
δn,l : (x1, . . . , xn) ↦→ (x1, . . . , x1,<br />
. . . , xn, . . . , xn)<br />
<br />
l1 times<br />
ln times<br />
2 Note that there are different definitions of a Wick monomial and polynomial in the<br />
literature. The definitions used here are not the most used ones.<br />
3 Remember that we considered the pullback of a distribution by a function in Sec-<br />
tion 4.2 of [4].<br />
72
Now by Theorem 4.5 from [4] we may define a Wick polynomial as an<br />
operator-valued distribution on D defined as the pullback of a Wick<br />
monomial by the a diagonal map δn,l. 4<br />
First we calculate the set of normals of the partial diagonal map.<br />
The transpose of the Jacobi matrix of the partial diagonal map is<br />
δ ′ n,l (x1, . . . , xn) t =<br />
⎛<br />
∂x1<br />
⎜<br />
⎝<br />
(δn,l)1<br />
⎞t<br />
. . . ∂xn(δn,l)1<br />
.<br />
. ..<br />
⎟<br />
. ⎠<br />
∂x1 (δn,l)n<br />
⎛<br />
1 . . . 1<br />
⎜ 0 . . . 0<br />
⎜<br />
= ⎜ .<br />
⎝ 0 . . .<br />
. . . ∂xn(δn,l)n<br />
0 . . .<br />
1 . . . 1 0 . . .<br />
0 1 . . . 1 0<br />
⎞<br />
0<br />
0 ⎟<br />
. ⎟ ,<br />
⎟<br />
0 ⎠<br />
0 . . . 0 1 . . . 1<br />
where row j contains lj ones. Hence<br />
δ ′ n,l (x1, . . . , xn) t η =<br />
⎛<br />
⎜<br />
⎝<br />
η1 + . . . + ηl1<br />
ηl1+1 + . . . + ηl1+l2<br />
.<br />
ηl1+...+ln−1+1 + . . . + ηl1+...+ln<br />
⎞<br />
⎟<br />
⎠ .<br />
Note that the matrix product is independent on x. The set of normals is<br />
Nδn,l = {(δn,l(x), η) ∈ ∆l(M) × R n |δ ′ n,l (x)t η = 0}<br />
= {(y, η) ∈ ∆l(M) × R n li+1 <br />
| ηli+j = 0 for all i = 0, . . . , n − 1},<br />
Where we have set l0 := 0.<br />
j=1<br />
Definition 9.5. By a Wick polynomial : φ l1 (x1) · · · φ ln (xn) : we mean the<br />
operator-valued distribution<br />
when it exists, that is, if<br />
.<br />
: φ l1 (x1) · · · φ ln (xn) : def<br />
= δn,l ∗ : φ(x1) · · · φ(xn) :<br />
Nδn,l ∩ WF (: φ(x1) · · · φ(xn) :) = ∅.<br />
4 Note that we proved the theorem exactly for the case of a diagonal map.<br />
73
Note that the definition agrees with Definition 9.3, since the restriction to<br />
the total diagonal ∆ n,(1,...,1)(M) simply returns the monomial itself.<br />
Translational invariance has been an important ingredient in the<br />
development of the method of Epstein-and Glaser. Unfortunately<br />
translations are not well-defined on general manifolds and have to be<br />
replaced by parallel transport. Therefore it is one of the key problems one<br />
has to face in the task of creating a microlocal formulation of the method<br />
(and thus make the method work on curved spaces) to find a suitable<br />
condition of smoothness to replace the translational invariance.<br />
The condition is related to graph theory. Letting G denote the set of all<br />
graphs, we define<br />
and<br />
Gn := {G ∈ G|G is non-oriented with vertices {1, . . . , n}}.<br />
E G := {e|e is an edge of G}.<br />
Further given e ∈ E G connecting i and j with i < j we say that i := s(e)<br />
and j := r(e) are the source and range of e, respectively.<br />
Definition 9.6. Let G ∈ Gn and M a manifold. An immersion (x, γ, k) is<br />
a triplet consisting of<br />
1. A map x mapping the vertices v of G to x(v) ∈ M<br />
2. a map γ mapping the edges e ∈ E G to null-geodesics γ(e) ⊂ M<br />
connecting x s(e) and x r(e)<br />
3. A map k mapping edges e ∈ E G to future directed covector fields<br />
k γ(e) := k(e), which are coparallel to the tangent vector ˙γe of the<br />
null-geodesic.<br />
Along with symmetry of Tn and causality we expected the S-matrix to<br />
have we have to add one more property:<br />
74
Microlocal Spectrum Condition.<br />
For the numerical distribution tn ∈ D ′ (M n ), n ≥ 2 of any time-ordered<br />
product,<br />
WF (tn) ⊂ Γn,<br />
where<br />
Γn =<br />
<br />
(x1, k1; . . . ; xn, kn) ∈ T ∗ M n \ {0}|∃g ∈ Gn and<br />
an immersion (x, γ, k) of G, in which ke is<br />
future directed whenever x s(e) /∈ J − (x r(e)) and<br />
such that ki = <br />
m:s(m)=i<br />
km(xi) − <br />
n:r(n)=i<br />
<br />
kn(xi) ,<br />
where t and s runs over all curves terminating and starting at xi,<br />
respectively.<br />
where<br />
J − (x) := {y ∈ M|y < x and there exists γ causal, connecting x and y}<br />
is the set of all points in M in the past of x that can be connected with x<br />
by a causal curve.<br />
The microlocal spectrum condition may be seen as a replacement of the<br />
spectrum condition (Axiom 4iii of the Wightman Axioms). Further it<br />
ensures that the wave front set has the following property.<br />
Figure 9.1: Illustration of Lemma 9.7.<br />
Lemma 9.7. Let (x, k) ∈ Γn then there exists a pair (xm, km) such that<br />
km /∈ V + (0).<br />
75
Proof. Let xm be a maximal point, that is xm /∈ J − (xi) for all i. Note that<br />
we may write<br />
km = −k1,m − · · · − km−1,m + km+1,m + · · · + kn,m,<br />
some of which might be zero, but not all since xm is connected to at least<br />
> 0 and<br />
one point. Now since xm is maximal, k0 1,m , . . . , k0 m−1,m<br />
k0 m+1,m , . . . k0 n,m < 0.<br />
The spectrum condition is important for the construction of the Scattering<br />
matrix, as the following theorem will witness.<br />
Theorem 9.8. If the spectral condition is satisfied, then<br />
Tn(x1, . . . , xn) : φ l1 (x1) · · · φ ln (xn) : (9.3)<br />
is a well-defined operator-valued distribution for any n and any choice of<br />
indices l1 . . . , ln on D. FIXME: dense invariant?<br />
Proof. Let Ψ ∈ D. Then by Definition 9.5<br />
: φ l1 (x1) · · · φ ln (xn) : Ψ<br />
is an operator-valued distribution. By Definition 9.2<br />
Thus FIXME: tjeck det følgede:<br />
WF (: φ(x1) · · · φ(xn) : Ψ) ⊂ (T ∗ M)−.<br />
WF (: φ l1 (x1) · · · φ ln (xn) : Ψ) = WF (δn,l ∗ : φ(x1) · · · φ(xn) : Ψ)<br />
= δn,l ∗ WF (: φ(x1) · · · φ(xn) : Ψ)<br />
⊂ M n × V ×n<br />
− ,<br />
where we have used that WF (f ∗ u) ⊂ f ∗ WF (u) by Theorem 4.2 in [4].<br />
Now remember our main result Theorem 4.10 of [4], that the product of<br />
two distributions u and v exists unless<br />
(x, k) ∈ WF (u) and (x, −k) ∈ WF (v)<br />
for some (x, k).<br />
But in our case the spectral condition insures that WF (Tn) ⊂ Γn and<br />
therefore by Lemma 9.7 there doesn’t exist a (x1, k1, . . . , xn, kn) ∈ Γn with<br />
ki ∈ V ×n<br />
+ for all i. Thus the product (9.3) is well-defined.<br />
76
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2nd edition Springer-Verlag, New York, 1995.<br />
[2] Gerd Grubb: Introduction to Distribution Theory - Lecture Notes,<br />
”http://www.math.ku.dk/∼grubb/distribution.htm”, Copenhagen,<br />
2003.<br />
[3] Romeo Brunetti and Klaus Fredenhagen: Microlocal Analysis and<br />
Interacting Quantum Field Theories: Renormalization on Physical<br />
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